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13700 | https://web.mit.edu/sahughes/www/8.022/lec16.pdf | Scott Hughes 12 April 2005 Massachusetts Institute of Technology Department of Physics 8.022 Spring 2005 Lecture 16: RL Circuits.
Undriven RLC circuits; phasor representation.
16.1 RL circuits Now that we know all about inductance, it’s time to start thinking about how to use it in circuits. A simple circuit that illustrates the major concepts of inductive circuits is this: V L R S2 S1 Suppose that at t = 0, we close the switch S1, leaving S2 open. How does the current evolve in the circuit after this?
Let’s first think about this physically. Lenz’s law tells us that the magnetic flux through the inductor does not “want” to change. Since this flux is initially zero, the inductor will impede any current that tries to flow it — the current will initially try to remain at zero.
As time passes, current will gradually leak through, and the magnetic flux will build up.
Eventually, it should saturate at I = V/R — at this point the most current possible is flowing through the circuit.
16.1.1 To Kirchhoffor not to Kirchhoff We’d now like to analyze this circuit quantitively as well. In all other circumstances in which we’ve examined a circuit, we’ve used Kirchhoff’s laws for this; for a a simple single-loop circuit, the relevant rule is Kirchhoff’s second law, “The sum of the EMFs and the voltage drops around a closed loops is zero”.
This rule is essentially a restatement of the old electrostatics rule that H ⃗ E ·d⃗ s = 0. This rule DOES NOT HOLD when we have a time changing magnetic field! In other words, inductors completely invalidate — at least formally — the foundation of Kirchhoff’s second law.
146 So what do we do? Let’s reconsider: in place of H ⃗ E · d⃗ s = 0, we have Faraday’s law: I ⃗ E · d⃗ s = −1 c dΦB dt = Eind = LdI dt .
The closed loop integral of ⃗ E gives us the induced EMF. I personally find it best to think of the inductor as an EMF source; it acts in the circuit effectively like a battery would. Even though Kirchhoff’s laws no longer stand, formally, on a very solid foundation, we can extend the formalism by thinking of inductors in the circuit as a source of EMF. The equations then carry over just fine; indeed, they carry over so well that almost all textbooks tell you to just apply Kirchhoff’s laws to circuits containing inductors without even mentioning this subtlety. It’s always worth bearing in mind that what you are really using is Faraday’s law.
I will try to refer to the loop rule as Faraday’s law; it’s such a habit to call it Kirchhoff’s law, though, that I will surely slip from time to time.
The tricky thing in these kinds of problems is getting the signs right. In all of these cases, Lenz’s law should guide you: the induced EMF will act in such a way as to oppose changes in the inductor’s magnetic flux. If a circuit initially has no current flowing through it, the inductor will oppose the build-up of current; if it initially does have current flowing, the inductor will try to keep that current flowing. Keep this intuition in mind and you should have no problem choosing the correct sign for the inductor’s EMF.
(Very detailed discussion of inductors and the failure of Kirchhoff’s laws can be found in Walter Lewin’s OCW 8.02 lectures from 2002; see Lectures 15 and 20 in particular.) 16.1.2 Back to our regularly scheduled program OK, let’s calculate. The EMF due to the battery is +V ; the EMF induced in the inductor is −LdI/dt. Why the minus sign? It’s going to fight the growth in current, and so must be opposite to the battery. Finally, there’s a voltage drop −IR at the resistor: V −IR −LdI dt = 0 .
Rearrange: dI dt = V −IR L = R L µV R −I ¶ .
Now divide; we end up with something that is of the form du/u: dI I −V/R = −R dt L so it is easy to integrate up, using the boundary conditions I = 0 at t = 0, I = I(t) at some general time t: ln "I(t) −V/R −V/R # = −Rt L or I(t) = V R ³ 1 −e−Rt/L´ .
147 As we guessed on physical grounds, the current builds up, with a “growth timescale” L/R.
The timescale L/R plays a role in LR circuits similiar to that played by the timescale RC in circuits with capacitors.
Suppose we let this circuit sit in this state for a very long time, so that the current saturates at its equilibrium value, I = V/R. What happens if we now simultaneously open S1 and close S2? Physically, our expectation is that Lenz’s law will try to hold the current flow constant — the inductor says “I don’t want the flux through me to change!” However, we know that the resistor dissipates energy, so the current must bleed away — we should see some kind of exponential decay.
Time to calculate. Going around the circuit, we have −LdI dt −IR = 0 .
The first term is the EMF due to the inductor, −LdI/dt. Why the minus sign? By Lenz’s law, the inductor always fights change. If the current “wants” to be reduced, the sign of the inductor’s EMF needs to fight that. The second term is the usual voltage drop across a resistor. Rearranging leads to the easily integrated equation dI I = −R dt L , which in turn leads to the exponential solution I(t) = I(t = 0)e−Rt/L .
(Note that we have redefined the meaning of t = 0 — this is now the time at which we open S1 and close S2.) When thinking about any circuit with inductance, the key thing you should keep in mind is that inductors cause currents to have a kind of “inertia”. If there is no current flowing, the inductor will force the current to build up gradually; if there is current flowing, the inductor will force it to change gradually. If current is already flowing and there is a sudden break in the circuit, the inductor will do whatever it takes to keep the current flowing. A consequence of this is that when we open switches or break circuits very suddenly, we may see a large spark! That spark is because of Lenz’s law — inductors in the circuit refuse to let the flux through them change. If that means that they must force the current to flow through the air, creating a huge, fiery spark — so be it.
16.1.3 Unit check Before moving on, let’s do the traditional sanity check on our units: the ratio L/R should have the dimension of time for any of the above results to make sense. In cgs, the units of L are sec2/cm; the units of R are sec/cm. So L/R is indeed a time measured in seconds.
In SI, the unit of L is the Henry, which is Volt-sec/Amp; the unit of R is the Ohm, which is Volt/Amp. So L/R is again a time measured in seconds. Three cheers for dimensional analysis.
16.2 LC circuits Exponential decay (or saturation) is dull. We’re now going to look at a circuit configura-tion that does something new. Consider what happens when we connect an inductor to a capacitor: 148 L C Suppose that the capacitor is charged up to a charge separation Q0. We close the switch at t = 0. How does the circuit behave after this?
As usual, our governing analysis of the system comes from Kirchhoff/Faraday: going around the circuit, we have Q C −LdI dt = 0 where Q is the charge on the capacitor, and I is the current flowing through the inductor.
To get anywhere with this, we of course need to relate the current flowing to the charge on the capacitor. Since the current comes from charge flowing offthe capacitor, we have I = −dQ dt .
We thus finally end up with the differential equation d2Q dt2 + 1 LC Q = 0 .
Differential equations of this form are best solved via an age old technique: guessing. In particular, guided by our experience with other systems that have a similar form, we guess that the solution has the form Q(t) = A cos ω0t + B sin ω0t .
We plug this solution back into our differential equation; the key thing we need is d2Q dt2 = −ω2 0A cos ω0t −ω2 0B sin ω0t = −ω2 0Q .
149 Our assumed solution works fine provided we set the angular frequency ω0 = 1 √ LC .
(You should be able to verify immediately that this has the correct units.) The constants A and B are free to be adjusted to match our initial conditions — the value of the charge on the capacitor Q at t = 0, and the current flowing in the circuit at t = 0: Q(0) = Q0 = A cos(0) + B sin(0) →A = Q0 .
I(0) = − ÃdQ dt !
t=0 = 0 = −ω0A sin(0) + ω0B cos(0) →B = 0 .
We now have the full solution to this problem: the charge on the capacitor, as a function of time, is Q(t) = Q0 cos ω0t .
Equivalently, the current flowing in the circuit, I(t) = −dQ/dt, is I(t) = ω0Q0 sin ω0t = Q0 √ LC sin ω0t .
This sinusoidal behavior is the generic feature of LC circuits. Notice that the charge has maximum magnitude when the current is zero, and vice versa — the current and the charge are “90◦out of phase” with each other.
A nice way to make sense of the current and charge behavior is by looking at the energetics of the circuit elements. The energy stored in the capacitor, as a function of time, is UC(t) = Q(t)2 2C = Q2 0 2C cos2 ω0t .
The energy stored in the inductor, as a function of time, is UL(t) = 1 2LI(t)2 = 1 2 1 LC LQ2 0 sin2 ω0t = Q2 0 2C sin2 ω0t .
150 The energy in the inductor and the energy in the capacitor each oscillate. Note that their sum is constant, and is exactly equal to the total energy that was initially stored in the capacitor: Utot = UC(t) + UL(t) = Q2 0 2C ³ cos2 ω0t + sin2 ω0t ´ = Q2 0 2C .
Conservation of energy works! In this case, there’s an even cooler interpretation. The energy in the capacitor is stored as electric fields; that in the inductor is stored as magnetic fields.
What we are seeing here is that the system’s total energy is just sloshing back and forth, from electric to magnetic to electric to ... This circuit is just an electromagnetic swing set!
16.3 RLC circuits The circuit we just discussed is an interesting one ... but it’s also somewhat fictional, since we are idealizing it as having no resistance. Since all circuits have some resistance1, the RLC circuit sketched here is far more relevant to real world applications than the LC circuit: C L R How would we expect this circuit to behave? Our intuition on the energetics of LC circuits is helpful here: qualitatively, we should expect the same “sloshing” of energy between the inductor and capacitor. However, we now expect that some energy will be dissipated through each cycle because of the circuit’s resistance. Our solution for the charge and the current in the circuit should look like a decaying oscillation of some kind.
To make progress, we need to do some calculating. We again assume that the capacitor initially has some charge Q0 on it, and that we close the switch at t = 0. Faraday/Kirchhoff then tells us how to set up the equations for this circuit: Q C −IR −LdI dt = 0 1Unless they are superconducting, but that opens a whole new barrel of worms ...
151 which reduces to d2Q dt2 + R L dQ dt + 1 LC Q = 0 .
The usual way to solve these equations is the tried and true friend of the physicist, inspired guessing: we assume that the solution takes the form Q(t) = e−t/τ (A cos ωt + B sin ωt) — very similar to the form we found for the LC circuit.
The major new feature is the exponentially decaying prefactor. We could then plug this assumed form of Q(t) into the differential equation and then grind out a lot2 of algebra. After the smoke clears, the algebra will tell you that, for this solution to work, τ and ω must take on certain values.
The constants A and B will be left as adjustable parameters to be used for matching initial conditions.
Feel free to pursue this route on your own!
We will instead use this moment as an opportunity to introduce a different way of thinking about this problem: we will use complex exponentials. For the remainder of this lecture, I will assume that you are familiar with the supplementary notes on complex numbers I put together (available under “Handouts” on the 8.022 webpage).
The key idea is to take the solution Q(t) to be the real part of a complex function ˜ Q(t): Q(t) = Re h ˜ Q(t) i .
Then, as long as the resistance, capacitance, and inductance are themselves real numbers, the complex charge ˜ Q is also a solution to the differential equation: d2 ˜ Q dt2 + R L d ˜ Q dt + 1 LC ˜ Q = 0 .
Our assumption is that ˜ Q = Aeiφ0eiαt .
As we will see shortly, A and φ0 (the “amplitude” and the “phase”) are the constants that we adjust to match initial conditions. We will now work out what α is. Bear in mind that α may itself be a complex number!!!
To proceed, we need to know the derivatives of ˜ Q. This is where the magic of using com-plex exponentials for solving differential equations comes in: This technique turns differential equations into algebraic ones. In particular, we have d ˜ Q dt = iα ˜ Q d2 ˜ Q dt2 = −α2 ˜ Q .
Plugging this back into the governing differential equation, we get ˜ Q µ −α2 + iαR L + 1 LC ¶ = 0 2A LOT.
152 which means that α must satisfy α2 −iαR L −1 LC = 0 .
This is just a quadratic equation, whose solution is α = i R 2L ± s 1 LC −R2 4L2 .
We thus have two solutions for ˜ Q: ˜ Q+ = Aeiφ0e−Rt/(2L)e+it√ 1/(LC)−R2/(4L2) ˜ Q− = Aeiφ0e−Rt/(2L)e−it√ 1/(LC)−R2/(4L2) .
The distinction between these two solutions is physically meaningless since they have the same real part! We might as well just keep one (the + solution, say) and be done with it.
Taking the real part, our final solution becomes Q(t) = Ae−Rt/(2L) cos (ωt + φ0) where ω = s 1 LC −R2 4L2 .
Although the use of complex numbers may seem kind of counterintuitive here, it makes solving this problem almost trivial. Rather than needing to fight through trig functions and their derivatives, you just need to solve one quadratic equation.
Complex numbers make for simple analysis.
16.3.1 The “weak damping” limit An important limit of the above circuit is that of weak damping.
This holds when the resistance is “small”, in a well-defined way.
To see what constitutes weak damping, we compute the current from our charge solution: I(t) = −dQ dt = Ae−Rt/(2L) · ω sin (ωt + φ0) + R 2L cos (ωt + φ0) ¸ .
The current consists of the sum of two different sinusoidal functions. The amplitude of this first term is proportional to the frequency ω; the amplitude of the second is proportional to the damping rate R/(2L).
Weak damping means that the frequency is much larger than the damping rate: ω ≫R 2L .
If this is the case, then the second term in the current may be ignored, to high accuracy, relative to the first term: I(t) ≃Ae−Rt/(2L)ω sin (ωt + φ0) .
153 We can actually make one further useful simplification in this limit: if ω ≫R/2L, then you should be able to convince yourself quite easily that ω ≡ s 1 LC −R2 4L2 ≃ s 1 LC = ω0 .
In other words, when we have weak damping, the oscillation frequency is — to high accuracy — the same frequency as we would have found if the resistor were not present!
In the weak damping limit, it is now simple to match the constants A and φ0 to our initial conditions: Q(t = 0) = Q0 = A cos(φ0) I(t = 0) = 0 = Aω0 sin(φ0) .
The solution to these equations is A = Q0 φ0 = 0 .
Thus, our “full” solution, in the weak damping limit, for this circuit is Q(t) = Q0e−Rt/(2L) cos ω0t I(t) = ω0Q0e−Rt/(2L) sin ω0t ω0 = s 1 LC .
A similar calculation can be done without assuming weak damping, but the details are much messier.
16.3.2 The quality factor In the weak damping limit, the energy content of the circuit behaves in a very nice way. The energy of the capacitor is given by UC(t) = Q(t)2 2C = Q2 0 2C e−Rt/L cos2 ω0t .
The energy of the inductor is given by UL(t) = 1 2LI(t)2 = 1 2ω2 0LQ2 0e−Rt/L sin2 ω0t = Q2 0 2C e−Rt/L sin2 ω0t .
154 (In going from the second to the third line, I’ve used ω2 0 = 1/LC.) The sum of these energies is Utot(t) = UC(t) + UL(t) = Q2 0 2C e−Rt/L ³ cos2 ω0t + sin2 ω0t ´ = Q2 0 2C e−Rt/L .
Since Q2 0/2C is the initial energy stored in the system, this means that the total energy content just exponentially damps away. (If we are not in the weak damping limit, a similar result holds, though the details are not quite as nice.) Scientists and engineers like to be able to quantify how “good” an RLC circuit is, and so we invented a quantity called the circuit’s quality Q (not to be confused with its charge). A circuit with high quality is one that oscillates many times before it loses a significant amount of energy. To make this somewhat vague notion precise, the following definition is typically adopted: “The quality factor Q is the number of radians through which a damped system oscillates by the time that its energy content has fallen by a factor of 1/e.” This may sound a little odd, but it deconstructs into a very simple mathematical state-ment. First, how much time does it take for the energy to fall by a factor of 1/e? By the above discussion, this is easy: T(energy to fall by 1/e) = L/R .
Plug this time into the energy stuffdiscussed above, and you see that the energy is at 1/e of its initial level.
How many radians does the system oscillate through in this time? We just multiply the time T by the system’s angular frequency ω ≃ω0! The quality factor of an RLC circuit is Q = ωL R ≃ ω0L R .
As we will see when we begin exploring AC circuits, there is another way to understand the quality factor of circuits, related to how a circuit behaves “on resonance” (and very close to resonance). We will define the quality factor in what seems to be a completely different way, and will end up with exactly the same result! This is not an accident of course, but is deeply connected to how energy is stored and dissipated in circuits.
The concepts we hit in this lecture — exponential decay of an oscillation; quality factor — and that we’ll hit in the next few lectures (e.g., resonance) pop up in an enormous number of contexts. For example, the quality factor of a distorted black hole — the number of radians worth of gravitational waves produced by the oscillations of its distorted event horizon emitted as it settles down to an equilibrium state — is given approximately by QBH ≃2 µ 1 −c G S M 2 ¶−9/20 (where G is the gravitational constant, c is the speed of light, S is the magnitude of the black hole’s spin angular momentum, and M is its mass). There are many other examples of quality factors that pop up repeatedly in science and engineering; you will see many of these notions again!
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Properties of Similar Polygons
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Similar Polygons
Similar polygons are shapes that have the same shape but not necessarily the same size. They maintain proportional side lengths and identical angles.
The following diagrams show the properties of similar polygons. Scroll down the page for more examples and solutions.
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Two polygons are similar if:
Corresponding angles are equal.
Corresponding sides are proportional.
Notation:
If polygon ABCD is similar to polygon WXYZ, we write:
ABCD ∼ WXYZ
Scale Factor
The scale factor is the ratio of any two corresponding sides in similar polygons.
If Polygon 1 has side lengths s 1, s 2, s 3, … and Polygon 2 has corresponding side lengths S 1, S 2, S 3, … then:
S 1 s 1=S 2 s 2=S 3 s 3=…=k S 1 s 1=S 2 s 2=S 3 s 3=…=k
where k is the scale factor.
If k>1, Polygon 2 is an enlargement of Polygon 1.
If 0<k<1, Polygon 2 is a reduction of Polygon 1.
If k=1, the polygons are congruent.
Ratio of Perimeters
The ratio of the perimeters of two similar polygons is equal to their linear scale factor.
Ratio of Areas
The ratio of the areas of two similar polygons is equal to the square of their linear scale factor.
Similar Polygons
Two polygons are similar if their corresponding angles are congruent and the corresponding sides have a constant ratio (in other words, if they are proportional).
Typically, problems with similar polygons ask for missing sides. To solve for a missing length, find two corresponding sides whose lengths are known. After we do this, we set the ratio equal to the ratio of the missing length and its corresponding side and solve for the variable.
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Similar Figures - Similar Polygons
Students learn that similar polygons have the same shape, and if two polygons are similar, then the corresponding angles are congruent, and the corresponding sides are in proportion. Students also learn that the scale factor is ratio of the lengths of two corresponding sides.
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13702 | https://andrewalpert.com/blog/the-steps-of-a-criminal-trial/ | Published Time: 2025-02-19T23:40:17+00:00
The Steps Of A Criminal Trial | Alpert Schreyer Criminal Defense Attorneys
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The Steps Of A Criminal Trial
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Alpert Schreyer Criminal Defense Attorneys | February 19, 2025 | Criminal Defense
The criminal court process can be long. Once you get to the trial phase, there are six main steps. Each step plays an important part in the jury trial process.
Pretrial Motions
Pretrial motions are requests a prosecutor or defense attorney makes to the court before trial. Some complex pretrial motions are filed weeks or months before the trial date, while other simpler matters are handled immediately before the trial begins.
Pretrial matters can include:
Motions to limit evidence
Motions to suppress because of an unlawful arrest, search or seizure
Motions to dismiss
Motion to close the courtroom
Motion to sequester witnesses
Pretrial motions are a helpful way that a lawyer can control how the trial will proceed.
Jury Selection
Jury selection is one of the most important parts of the trial process since the jury renders the verdict. During jury selection, the judge brings in the jury venire. This is a larger panel of jurors that the attorneys question to get to the final jurors.
In Maryland, there are 12 jurors in a felony trial, plus one or two alternate jurors. Each side and the judge can ask the venire questions to eliminate jurors who have biases. If someone has a bias and is unable to be an impartial juror, they will be struck for cause. There are unlimited strikes for cause.
Additionally, each side has four peremptory strikes in most cases. A peremptory strike means that the lawyer can exclude a juror for any reason aside from racial bias.
Opening Statements
After the jury is sworn in, both sides will proceed to opening statements. During opening statements, the lawyers will lay out their case and give the jurors a roadmap of what to expect. The opening statement is not an argument but rather an overview of the law and the expected evidence.
Witness Testimony
After opening statements, the prosecution will call their witnesses and enter relevant evidence. The prosecutors will question their witnesses during direct examination, and then the defense will cross-examine the witnesses. Cross-examination is a helpful way for the defense to poke holes in the prosecutor’s evidence.
After the prosecution has finished calling witnesses, they will “rest.” Then, the defense has an opportunity to call witnesses. However, the defense is not obligated to call any witnesses or enter any evidence because it does not have a burden of proof. Instead, the prosecution has the burden to prove the case beyond a reasonable doubt.
If the defense does call witnesses, the prosecution can cross-examine the witnesses just as the defense had an opportunity to cross-examine earlier.
Closing Arguments
After all of the evidence is presented, both sides will make closing arguments. During closing arguments, the lawyers will discuss the elements of the crime and the evidence presented during the trial and argue for either a guilty or not guilty verdict.
Unlike during opening statements, closing arguments tend to be more passionate. The lawyers have already laid out the evidence for or against the case and are making a final plea to the jurors.
Jury Deliberation
After closing arguments, the judge will read jury instructions to the jurors. The jury instructions are instructions that the jurors must follow while deliberating. Jury instructions include instructions about:
The law
The burden of proof
The elements of the crime
The jury deliberation process
Impartiality
Excluded evidence
When they enter the jury deliberation room, the jurors will receive copies of these instructions and all of the physical evidence.
Deliberating can take anywhere from hours, days, or even weeks in complicated cases. If the jury has questions about the evidence or law during deliberation, they can write down the question and submit it to the judge. The judge will answer the question with both the prosecutor and defense attorneys present.
Once the jury reaches a unanimous verdict, it will submit a jury verdict form to the judge. If the jury cannot reach a unanimous verdict, it results in a hung jury. The judge may then declare a mistrial, though further deliberation may be allowed before that decision is made. Usually, the foreperson or court clerk reads the verdict in the courtroom. If the verdict is guilty, the case will proceed to sentencing. The defendant is usually free to go if the verdict is not guilty.
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13703 | https://www.quora.com/What-is-Arithmetico-Geometric-series | Something went wrong. Wait a moment and try again.
Geometric Sequences
Finite Series
Arithmetic-Geometric Mean
Number Series Mathematics
Algebra Series
Series of Numbers
Numerical Series
5
What is Arithmetico-Geometric series?
Sonu Jha
Former Faculty at Resonance Eduventures (2016–2022)
·
5y
All details about
ARITHMETICO GEOMETRIC SERIES
is……
I hope it helps you.
All details about
ARITHMETICO GEOMETRIC SERIES
is……
I hope it helps you.
Related questions
How do I find N in a geometric series?
What is the geometric series 1+2+4+…2⁴⁸+2⁵⁰?
What is the infinite geometric series of 0.7?
What is the closed form of a geometric series?
How can I apply for a geometric series?
Lukas Schmidinger
I have graduate CS and my studies included math courses.
·
Author has 27.7K answers and 14.9M answer views
·
1y
Originally Answered: In a series, what is the difference between an arithmetic series and a geometric series?
·
In a series, what is the difference between an arithmetic series and a geometric series?
An arithmetic series is the sum of a arithmetic progression. An arithmetic progression has a common difference, the general nth element is an=a1+(n−1)d with a1 being the first element and d as the difference. An example can be the positive integers: an=n (a1=1 and d=1) gives 1, 2, 3, …. Thus the sum is sn=n∑k=1ak=na1+n2−n2⋅d (note that in our example this works out great to be the formula for triangular numbers n2+n2).
A geometric series is the sum of
In a series, what is the difference between an arithmetic series and a geometric series?
An arithmetic series is the sum of a arithmetic progression. An arithmetic progression has a common difference, the general nth element is an=a1+(n−1)d with a1 being the first element and d as the difference. An example can be the positive integers: an=n (a1=1 and d=1) gives 1, 2, 3, …. Thus the sum is sn=n∑k=1ak=na1+n2−n2⋅d (note that in our example this works out great to be the formula for triangular numbers n2+n2).
A geometric series is the sum of a geometric progression. A geometric progression has a common ratio, the general nth element is an=rn−1a1 with a1 once again being the first element and r being the common ratio. An example might be the powers of two: an=2n−1 (a1=1 and r=2) gives 1, 2, 4, …. So the sum is sn=n∑k=1ak=a1(rn−1)r−1 (for r=1 we have an=a1 and thus also sn=na1).
Manjunath Subramanya Iyer
I am a retired bank officer teaching maths
·
Author has 7.2K answers and 10.4M answer views
·
Updated 6y
A series of the form
a + (a+d)r + (a+2d)r^2 + (a+3d)r^3+……….
is called Arithmetico-Geometric series.
Notice that it has the ‘flavour’ of both AP (due to the presence of common difference d) as also GP(due to the presence of the common ratio r).
Assistant
Bot
·
1y
An Arithmetico-Geometric series is a series that combines the characteristics of both arithmetic and geometric sequences. It is formed by multiplying each term of an arithmetic sequence by a corresponding term of a geometric sequence.
General Form
The general form of an Arithmetico-Geometric series can be expressed as:
S=a1r0+(a1+d)r1+(a1+2d)r2+(a1+3d)r3+…+(a1+(n−1)d)rn−1
Where:
- a1 is the first term of the arithmetic sequence.
- d is the common difference of the arithmetic sequence.
- r is the common ratio of the geometric sequence.
- n is the number of
An Arithmetico-Geometric series is a series that combines the characteristics of both arithmetic and geometric sequences. It is formed by multiplying each term of an arithmetic sequence by a corresponding term of a geometric sequence.
General Form
The general form of an Arithmetico-Geometric series can be expressed as:
S=a1r0+(a1+d)r1+(a1+2d)r2+(a1+3d)r3+…+(a1+(n−1)d)rn−1
Where:
- a1 is the first term of the arithmetic sequence.
- d is the common difference of the arithmetic sequence.
- r is the common ratio of the geometric sequence.
- n is the number of terms in the series.
Example
For example, consider the series:
S=2⋅1+(2+1)⋅2+(2+2)⋅4+(2+3)⋅8
This can be expanded as:
S=2+6+16+40
Sum of an Arithmetico-Geometric Series
The sum of an Arithmetico-Geometric series does not have a simple closed form like a standard arithmetic or geometric series. However, it can often be evaluated using techniques such as manipulation of the series, differentiation, or generating functions.
Applications
Arithmetico-Geometric series appear in various areas of mathematics, including calculus, numerical methods, and in solving certain types of recurrence relations. They can be particularly useful in problems involving sequences and series where both additive and multiplicative behaviors are present.
If you have a specific series in mind or need help with a particular problem related to Arithmetico-Geometric series, feel free to ask!
Related questions
What is the difference between a geometric sum and a geometric series?
Could someone explain to me why a geometric series is called geometric?
What is a non-converging geometric series? What is the mathematical definition of a non-converging geometric series (with example)?
How do you write the following as a geometric series, 0.090 909 09?
The sum of the first n terms of an AP is 5n2−2n. What is its 50th term?
Michael Paglia
Author has 33.3K answers and 5.3M answer views
·
2y
Originally Answered: What is the difference between an arithmetic series and a geometrical series?
·
Arithmetic is having a number whose next number and prededing number is equal
Same number added each time
2..5…8..11. Add 3..
Geometric a little tougher to explain
A Fourier series is Geometric
That decreases by a fraction
Infinte
Ill say cutting something in half…again and again
Half of 8 is 4 difference is 4
Half of 4 is 2 ..difference is 2
Not 4…. the difference of the first one
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Ron Auerbach
Job search expert, author, career coach, & resume writer
·
Author has 13.4K answers and 20.5M answer views
·
1y
Originally Answered: In a series, what is the difference between an arithmetic series and a geometric series?
·
Arithmetic
series = ones where the each number changes by the same amount. For example: 2, 4, 6, 8. Here, you’re just adding 2 to the previous number to get the next. By contrast, a
geometric
series is one where the the difference isn’t the same but the ratio of the numbers is. For example, 2, 4, 8, 16. With this series, the difference between the various numbers isn’t always the same as with an arithmetic series. But you’re multiplying each number by 2. So the ratio of one number to the next remains the same.
Sanchit Hajela
Student
·
5y
In mathematics, an arithmetico–geometric sequence is the result of the term-by-term multiplication of a geometric progression with the corresponding terms of an arithmetic progression . Have a look for broad perspective !
What is Arithmetico–Geometric Sequence? - A Plus Topper
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George Ivey
Former Math Professor at Gallaudet University
·
Author has 23.7K answers and 2.6M answer views
·
1y
Originally Answered: In a series, what is the difference between an arithmetic series and a geometric series?
·
In an arithmetic series, the difference between two successive terms is a constant: an=a+d(n−1) where a is the first term d is the constant difference. 2+ 5+ 8+ 11+ 14+17+… is an arithmetic series with a= 2 and d= 3,
In a geometric series, the ratio of two successive terms is a constant: an=arn−1 where a is the first term and r is the constant ratio. 2+ 6+ 18+ 54+ 162+ … is a geometric series with a= 2 and r= 3.
Melissa Dalis
Majored in CS & Math at Duke
·
Upvoted by
Justin Rising
, PhD in statistics · Author has 174 answers and 1.9M answer views
·
11y
Related
How do you sum a series that is not a geometric series?
We want to find the inifinite sum of n2n=n(1/2)n
We can start with the Power Series, as follows:
11−x=1+x+x2+...+xn+...
Differentiating both sides, we get:
1(1−x)2=1+2x+3x2+...+nxn−1+...
x(1−x)2=x+2x2+3x3+...+nxn+...
Notice that the nxn term is what we were trying to find with n(1/2)n. So plug in x=1/2 to get:
1/2(1−1/2)2=1/21/4=2
We want to find the inifinite sum of n2n=n(1/2)n
We can start with the Power Series, as follows:
11−x=1+x+x2+...+xn+...
Differentiating both sides, we get:
1(1−x)2=1+2x+3x2+...+nxn−1+...
x(1−x)2=x+2x2+3x3+...+nxn+...
Notice that the nxn term is what we were trying to find with n(1/2)n. So plug in x=1/2 to get:
1/2(1−1/2)2=1/21/4=2
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Roman Andronov
Solving technical problems is my day job.
·
Author has 427 answers and 13.4M answer views
·
Updated 10mo
Related
Why is a geometric series useful? What are its applications?
Geometric series is useful because it can be used as a model of real life situations which can find its application in physics, for example.
Consider the following
Problem
From a given height h0 we, strictly horizontally, launched with a given initial velocity v0 a material point (assume a small spherical chicken of uniform density).
It is known that the point bounces off of a hard perfectly horizontal surface in such a way that after each bounce its vertical velocity diminishes at a constant rate r - that is the ratio of the point’s vertical velocity after each bounce to its vertical velocity b
Geometric series is useful because it can be used as a model of real life situations which can find its application in physics, for example.
Consider the following
Problem
From a given height h0 we, strictly horizontally, launched with a given initial velocity v0 a material point (assume a small spherical chicken of uniform density).
It is known that the point bounces off of a hard perfectly horizontal surface in such a way that after each bounce its vertical velocity diminishes at a constant rate r - that is the ratio of the point’s vertical velocity after each bounce to its vertical velocity before that bounce is exactly 0<r<1.
Question:
assuming the absence of surface friction, air drag and resistance and wind, under uniform gravity, v0<<c and so on (under ideal, high school) conditions - what will the horizontal displacement of the point be after it stops bouncing?
Solution
Under the mentioned ideal conditions the point’s horizontal velocity will always remain constant - and equal to v0.
On the first touchdown, which, loosely speaking, represents half a bounce or half a parabola, our point will be x0 units away from origin and after that it will make full bounces or trace full parabolas and thus it will cover xk units of distance between consecutive bounces starting with k=1 as shown below:
Thus, if we designate the point’s horizontal displacement as X, we shall have:
X=x0+n∑k=1xk(1)
for some, unspecified for now, number of bounces n. We may think of (1) as a formula that captures a certain steady, bouncy, state.
Prior to the first touchdown, our point is free-falling from the initial height h0. Hence, for the initial flight time t0 we have:
h0=gt202(2)
from where:
t0=√2h0g(3)
Therefore, for x0 we have:
x0=v0t0=v0√2h0g(4)
We now will be working with the vertical component of the point’s velocity so let us designate it as vy.
From the law of conservation of energy it follows that the point will touch the ground initially with vy0 such that:
vy0=√2gh0(5)
Consequently, after that bounce it will get airborne with the (initial) vertical speed vy1 which is known (as per the problem statement):
vy1=vy0⋅r=r√2gh0(6)
After the second bounce the point will get airborne with the (initial) vertical speed vy2 such that (using vy1 from (6)):
vy2=vy1⋅r=vy0⋅r⋅r=r2√2gh0(7)
After the third bounce the point will get airborne with the (initial) vertical speed vy3 such that (using vy2 from (7)):
vy3=vy2⋅r=vy1⋅r⋅r=vy0⋅r⋅r⋅r=r3√2gh0(8)
and so on.
We thus argue that after the k-th bounce the point’s initial vertical velocity vyk will be:
vyk=rk√2gh0,k∈N(9)
For brevity we skip the deduction of the magnitude of the maximum height hk to which the point climbs under the given conditions and take it to be already known:
hk=v2yk2g(10)
Put vyk from (9) into (10):
hk=r2k2gh02g=h0r2k(11)
(Observe that the parabola’s peaks diminish faster than the vertical velocities do)
Consequently, the time tk to climb to the maximum height hk is:
tk=√2hkg(12)
Put hk from (11) into (12):
tk=√2h0r2kg=rk√2h0g(13)
But we take it that the time that it takes to climb to hk is equal to the time that it takes to descend from hk. Therefore, the total time Tk that the point stays airborne between two consecutive full bounces is simply twice that of tk:
Tk=2tk=2rk√2h0g(14)
And since the horizontal component of the point’s velocity, v0, always remains the same, it follows that for xk we have:
xk=v0⋅Tk=2v0rk√2h0g(15)
Putting (4) and (15) back into (1), we obtain:
X=v0√2h0g+2v0√2h0gn∑k=1rk(16)
or:
X=v0√2h0g(1+2n∑k=1rk)(17)
Now, how do we translate the soft stops bouncing into hard math?
One way - our point stops bouncing if after the n-th touchdown it fails to get airborne.
That, in turn, means that:
hn=h0r2n=0(18)
But that is mathematical heresy because h0 is a non zero real number, so is r and n is a non zero natural number.
When will (18) become a good mathematical citizen?
When we recall the concept of a limit:
limn→+∞h0r2n=0(19)
because, don’t forget, r<1. In other words, we tend n to (positive) infinity.
Thus one geometric series, with a common ratio r, comes to life:
X=v0√2h0g(1+2+∞∑k=1rk)(20)
But we already know that:
+∞∑k=1rk=r1−r(21)
Therefore, putting (21) into (20) we arrive at the answer:
X=v0√2h0g(1+2r1−r)=
v0√2h0g⋅1−r+2r1−r
or:
X=v0√2h0g⋅1+r1−r(22)
Is that it?
No. Why not? Because our spherical chicken of uniform density participates in a frequent flier miles program (who said that chickens can’t fly?)
Since we have not one but two geometric series going on here, we might as well enquire - how much vertical distance did our chicken cover before it got grounded?
Well, that one is easy:
H=h0+2+∞∑k=1hk=
(as per (11))
h0+2h0+∞∑k=1r2k=
h0(1+2+∞∑k=1r2k)=
h0(1+2r21−r2)=h0⋅1−r2+2r21−r2
or:
H=h0⋅1+r21−r2(23)
Is that it?
We have yet another geometric series here which was hiding in plain sight all along - the total time T that the chicken spent in the air:
T=t0+2+∞∑k=1tk=
(as per (3) and (13))
t0+2√2h0g+∞∑k=1rk=
√2h0g(1+2+∞∑k=1rk)=
√2h0g(1+2r1−r)=√2h0g⋅1−r+2r1−r
or:
T=√2h0g⋅1+r1−r(24)
Extra for experts: you can now embellish the problem with additional variances: let v0 point up (or down) with respect to the horizon by a non zero angle α, let the hard surface be inclined to the horizon under a non zero angle θ and so on.
In case this is of interest to anyone: I made the diagram in GeoGebra, taking:
r=0.7
h0=7
and normalizing the equation of the very first parabola:
y(x)=h0−g2v20⋅x2
in such a way that the leading coefficient of x2 is a unity:
g2v20=1
which essentially makes it:
y(x)=7−x2
The equations of the remaining two parabolas are left as an exercise for the reader to find.
For more information on and ideas about problem-solving in mathematics, physics and computer science please visit my YouTube channel ProbLemma.
Qiaochu Yuan
ex-PhD student in math
·
Author has 311 answers and 5.6M answer views
·
11y
Related
Why are geometric series named as "geometric" literally?
Daniel Melchor
IB Math Higher Level student, 2nd year
·
5y
Related
What is a geometric sequence in math?
A geometric sequence is a sequence of numbers where the common difference between each of them is a multiplication or division.
Let’s say for example that we have this sequence:
To find out what the common factor is we can simply divide two consecutive numbers (the second by the first one) so, for example, we could find out that the common factor in this series is 2 by dividing 64/32.
We could also have a sequence that always gets smaller instead of bigger. Let’s say, for example, that we have a sequence where the common factor is less than one. So 0<r<1:
To fin
A geometric sequence is a sequence of numbers where the common difference between each of them is a multiplication or division.
Let’s say for example that we have this sequence:
1,2,4,8,16,32,64…
To find out what the common factor is we can simply divide two consecutive numbers (the second by the first one) so, for example, we could find out that the common factor in this series is 2 by dividing 64/32.
We could also have a sequence that always gets smaller instead of bigger. Let’s say, for example, that we have a sequence where the common factor is less than one. So 0<r<1:
1,1/2,1/4,1/8,1/16…
To find out the common factor we can do the same thing and we get 1/2.
Some formulas we can use are the following:
To find any term (n) in the sequence we can use an=a1∗rn−1. So for example, if we wanted to find the 3rd term in the first example, we could multiply the first term times the common factor to the power of 3–1. So we would have
If a sequence has a common factor where [math]-1<r<1[/math], we can use the formula [math]\frac{a_1}{1-r}[/math] to find out what the sum of the infinite terms in the sequence is. For example, in the second example where [math]-1<\frac{1}{2}<1[/math], the infinite sum would be [math]\frac {1}{1–0.5}=2[/math]
Finally, if we want to sum only X terms, we can use the formula [math]\frac{a_1(r^{n}-1)}{r-1}[/math]. Where the sum of the first 10 numbers in the first example would be: [math]\frac {1(2^{10}–1)}{1}=1023[/math]
Edit: equations
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What is a geometric series?
A geometric series is one which has a constant ratio between successive terms. In a geometric series you always have an ‘r’ (common ratio) and an ‘a’ (first term).
The common ratio and first term are constant, they never change. Generally you will see a geometric series written:
[math]a + ar + ar^2 + ar^3 + ...[/math]
As you can see, the common ratio’s power equals the number term of the series. Let’s look at an example
let a = 1 and r = 2
[math]1 + 12 + 12^2 + 12^3 + ...[/math]
[math]1 + 2 + 4 + 8 + ...[/math]
This is a geometric series. Hope that helps!
Related questions
How do I find N in a geometric series?
What is the geometric series 1+2+4+…2⁴⁸+2⁵⁰?
What is the infinite geometric series of 0.7?
What is the closed form of a geometric series?
How can I apply for a geometric series?
What is the difference between a geometric sum and a geometric series?
Could someone explain to me why a geometric series is called geometric?
What is a non-converging geometric series? What is the mathematical definition of a non-converging geometric series (with example)?
How do you write the following as a geometric series, 0.090 909 09?
The sum of the first n terms of an AP is 5n2−2n. What is its 50th term?
What is the value of R of the geometric series?
What is the definition of a geometric series? What are some examples of when a geometric series converges but its sum diverges?
What is the value of A1 of the geometric series?
What is a hypergeometric series, and how is it related to a standard geometric series?
What is a hyper-geometric series? Give an example.
Related questions
How do I find N in a geometric series?
What is the geometric series 1+2+4+…2⁴⁸+2⁵⁰?
What is the infinite geometric series of 0.7?
What is the closed form of a geometric series?
How can I apply for a geometric series?
What is the difference between a geometric sum and a geometric series?
Could someone explain to me why a geometric series is called geometric?
What is a non-converging geometric series? What is the mathematical definition of a non-converging geometric series (with example)?
How do you write the following as a geometric series, 0.090 909 09?
The sum of the first n terms of an AP is 5n2−2n. What is its 50th term?
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13704 | https://www.utas.edu.au/mathematics-pathways/pathway-to-engineering/compulsory-modules | COMPULSORY MODULES, 1-3 - Mathematics Pathways | University of Tasmania
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COMPULSORY MODULES, 1-3
UTAS Home››Mathematics Pathways›Pathway to Engineering›COMPULSORY MODULES, 1-3
These modules are compulsory to receive RPL for the Mathematics Units as part of the Engineering Pathway for VET students.
Please complete all the exercises in each of the following module, including the final Quiz which are required for assessment.
Complex Numbers
Calculus:1 st and 2 nd order Differential Equations
Calculus:Partial Differentiation
Instructions:
1.Move through the tabs sequentially within each compulsory module. Ensure that you complete all the exercises prior to attempting the quiz at the end.
2.PLEASE NOTE: the quiz is assessed and required as evidence that you satisfactorily completed the requirements to receive RPL.
3.You will have the option to have a copy of your completed quiz emailed to you.Ensure that you save this copy as evidence that you have completed the module to receive RPL at your TAFE.
Complex Numbers
Calculus, 1st and 2nd Order Differential Equations
Calculus, Partial Differentiation
Complex numbers allow us to solve equations that have no ‘real solution’ and extend the number line into a 2 dimensional complex plane as shown below.
“think of Adam and Eve like an imaginary number, like the square root of minus one: you can never see any concrete proof that it exists, but if you include it in your equations, you can calculate all manner of things that couldn't be imagined without it.” Philip Pullman, The Golden Compass
Photo: The complex plane showing operations, and the Mandelbrot set.
Image by Kan8eDie (Own work) via Wikimedia Commons
Information
Definition
Application
Lessons
Ex. 1
Ex. 2
Review
Quiz
Course Information
Complex numbers are part of the national engineering unit MEM23004A Apply technical mathematics more details on unit content from training.gov.au
Required skill: to solve problems involving complex number quantities using the properties, operations and theorems of complex numbers.
There are no prerequisites required for this module; however, to complete the lessons and quiz below you should have an understanding of algebra and polynomial equations.
The Definition of a Complex Number
A complex numbers is a number that can be expressed in the form z = a + b i, where a and b are real numbers and i is the imaginary unit, satisfying i 2 = -1
For Example, -3.5 + 2 i is a complex number
Engineering Applications of Complex Numbers
Complex numbers are used by Electrical & Electronic Engineers to define the Alternating Current or AC concept of Impedance, and in Fourier analysis they are used in the processing of radio, telephone and video signals, see thispage for more details.
Mechanical & Structural Engineers use complex numbers to analyse the vibration of structures in machines, buildings and bridges, the behaviour of fluid flow around aircraft, and that of wind around buildings and bridges, preventing failures such as the Tacoma Narrows Bridge (Please watch the following video).
Complex numbers are also the basis for the Mandelbrot set, an example of 2D fractal geometry as shown in the image below.
Image from Wikipedia licensed under the Creative Commons Attribution 2.5 Generic Licence
For more details on the Mandelbrot Set, watch this TED Talk by Benoit Mandelbrot (17min).
Examples of approximate fractals in the natural world include leaves, mountain ranges, clouds, broccoli and many others; see this page for more examples.
Lessons on Complex Numbers
To get a better understanding of complex numbers review these lessons and quizzes from the following links:
Complex Numbers- Interactive Mathematics
Complex Numbers - Khan Academy
Watch the following video from Sixty Symbols on Imaginary Numbers
Basic Operations with Complex Numbers- Interactive Mathematics
Adding Complex Numbers- Khan Academy
Now go to the Interactive graphical addition, subtraction tool from geogebratube.org
Divide complex numbers -Khan Academy
Graphical Representation of Complex Numbers- Interactive Mathematics
Now take the interactive graphing quiz from Khan Academy to practice graphing complex numbers
Go to the Interactive graphical tool cartesian and polar coordinates from geogebratube.org to practice using the Argand Diagram and in exponential form
Solving quadratic equations with complex solutions - Interactive Mathematics
Use the interactive graphical tool from nrich.maths.org for practicing these ideas
All Khan Academy content is available for free at
Engineering Worked Examples and Practice Exercises
Example
Mechanical/Structural Engineering - Vibration
Machines and structures, such as a bridges or buildings are often subjected to varying loads or forces, which will cause the movement, or displacement, of a point on the machine or structure possibly causing vibration.
A weight and undamped spring are an example of simple harmonic motion
Image licensed under the Creative Commons Attribution-Share Alike 3.0 Unported license.
Engineers may also need to create forced vibration to move bulk materials such as fine powders, gravel and rocks in industries such as mining, food processing and cement manufacturing see examples here.
Common method to create forced vibration by the rotation of an offset mass.
Image licensed under the Creative Commons Attribution-Share Alike 3.0 Unported license from Brews ohare
Definitions: Motion of a point on a structure is given by displacement
x = XSinωt
Where, X= amplitude of single peak hence x varies between + X and - X,ω= frequency in radians/second,t = time in seconds.
Graphical representation of vibration by sine or cosine functions
Public domain image from the Wikimedia Commons
Differentiation of displacement with respect to time gives velocity, ν =ωΧCosωt
Differentiation of velocity with respect to time gives acceleration, a = -ω 2 ΧSinωt
These equations can be expressed in complex exponential form which simplifies many operations, especially when we have vibration in multiple planes in 3D space.
3D representation of the complex and real planes showing a complex function in exponential form e jz resulting from the addition of a sine and cosine function as described below.
Image licensed from Qniemiec under the Creative Commons Attribution-Share Alike 3.0 Unported license.
Using Euler’s Equation, e iϑ = Cosϑ + iSinϑ where ϑ= angle in radians, e= base of the natural logarithm or "Euler's number" approximately = 2.71828.
If we multiply Euler’s Equation by Χ and substitute ϑ =ωt we get
Χe iωt =ΧCosωt + iΧSinωt
Note similarity to our equations for displacement and velocity above.
When we plot this equation on the complex plane or Argand diagram we get a vector of length Χ rotating at ω radians/sec.
Where displacement x = Real part of Χe iωt and velocity v = Imaginary part of Χe iωt
Image from Smith, J.O. Introduction to Digital Filters with Audio Applications, online book, accessed 10/01/14, copy permission
This image is licensed under the Creative Commons Attribution-Share Alike 3.0 Unported license from Gonfer
Problem:
For the design of a vibrating feeder we need to determine the forces in the system. We can approximate the motion of the feeder to simple harmonic motion represented by the sine or cosine function.
Our vibrating feeder design has a single peak amplitude Χ= 2mm or 0.002m at a frequency ω =18 cycles/sec. Calculate the single peak velocity and acceleration.
Worked Solution:
We need to convert ω from cycles/sec to radians/sec, 1 cycle or revolution = 2 π radians, ω=18 x 2 π
ω=36 π radians/sec
Velocity,ν =ωΧCosωt, and peak velocity occurs when Cos ωt=1 (see graph above)
ν =ωΧ= 36 π 0.002 = 0.23 m/sec
Acceleration,a = -ω 2 ΧSinωt, and peak acceleration occurs when Sinωt = 1
a = -ω 2 Χ= (36 π)2 0.002 = 25.58 m/ sec 2
Practice exercise:
If a new bridge design has a single peak amplitude Χ= 3mm or 0.003m at a frequency ω= 5 cycles/sec find the single peak velocity and peak acceleration?
Check your answer to the practice exercise here.
Engineering Worked Examples and Practice Exercise
Example
Electrical/Electronic Engineering - Circuit Design
Electrical current measurements are important to people who want to use electricity in a creative way. Electrical current is used as a direct current (DC), the one directional flow of electric charge or alternating current (AC), where the flow of electric charge periodically reverses direction, see diagram below. When working with AC, complex numbers are needed.
Public domain image from the Wikimedia Commons
Electric current equations:
V= I x Z or Z = V/I
where, V = voltage, I = current and Z = impedance
Definitions:
Voltage: the difference in electrical charge between two points in a circuit, unit in volts, V
Current: the flow of electric charge, units are amperes or amps, A
Impedance: a measure of the opposition that a circuit presents to a current when voltage is applied.
The magnitude of complex impedance gives the ratio of the voltage amplitude to the current amplitude. In polar form Z = |Z|e iθ where θ =phase angle or difference between voltage and current, using complex notation Z = R + i X where, X = reactance and R = resistance, in DC circuits we only consider R units in ohms Ω.
In many applications the relative phase of the voltage and current is not critical so only the magnitude of the impedance is used or |Z|
Note: i is the imaginary unit, in electrical engineering j is used in place of i to avoid confusion with the symbol for electric current = i
Public domain image from the Wikimedia Commons
Series circuit diagram
Problems:
The impedance in one part of a circuit shown above is 4 + 12 i ohms, the impedance in another part of the circuit connected in series is 3 - 7 i ohms. What is the total impedance in the circuit?
ZT = Z1 + Z2... = (4 + 12 i ) + (3 - 7 i )
ZT = 7 + 5 i ohms
The voltage in a circuit is 45 + 10 i volts and the impedance is 3 + 4 i ohms, what is the current?
If the current in a circuit is 8 + 3 i amps and the impedance is 1 - 4 i ohms, what is the voltage?
I = 8 + 3 i amps
Z = 1 - 4 i ohms
V = I x Z
V = (8 + 3 i ) x (1 - 4 i )
V = (8 -32 i + 3 i - 12 i 2 ) = (8 - 32 i + 3 i + 12)
V = 20 - 29 i volts
Practice exercises:
The voltage in a series circuit is 5 +10 i volts and the impedance is 1 + 2 i ohms, what is the current?
If the current in a series circuit is 9 - 4 i amps and the impedance is 3 + 4 i ohms, what is the voltage?
Note: For components connected in parallel, the voltage across each circuit element is the same; the ratio of currents through any two elements is the inverse ratio of their impedances.
Public domain image from the Wikimedia Commons
Parallel Circuit Diagram
Check your answers to the practice exercises here.
Review Exercises
Basic Interactions with Complex Numbers from Interactive Mathematics with solutions
Complex Numbers review from SOS Mathematics with solutions
Reference Web Sites
These are some useful websites with information on Complex Numbers
Interactive Mathematics uses various math applets to enhance mathematics lessons.
Topics range from grade 8 algebra to college-level Forier and Laplace Transformations.
History of complex numbers.
Complex Numbers Quiz
When you have completed the lessons, verify your knowledge with this multiple choice Complex Numbers Quiz
If you need to repeat the quiz, click here(please note: you will not be able to repeat the first quiz)
PLEASE NOTE:
This Quiz is assessed and required as evidence that you have satisfactorily completed the requirements to receive RPL at your TAFE
Please attempt all questions on the quiz
To pass this quiz you will need to achieve at least 70%
You will need to enter your name and email address in the quiz to have quiz results emailed to you, ensure that you save and print out these results as evidence that you have passed the module, your TAFE will ask to see these results.
Calculus is one of the greatest achievements of the human mind separately invented by Newton and Leibniz. Calculus allows us to understand things that change, originally to understand motion, or the change in position over time. Calculus explains a multitude of real world phenomenon including mechanics and electromagnetism in engineering.
Photo: Solution to the differential equations for heat transfer in a pump casing model using finite element modelling software.
This file is licensed under the Creative Commons Attribution-Share Alike 3.0 Unported license
Information
Definition
Applications
Lessons
Review
Examples
Quiz
Course Information
Differential equations are part of the national engineering unit MEM23007A Apply calculus to engineering tasks. More details on unit content from training.gov.au
Required skill: identifying and solving simple first and second order differential equations.
The prerequisite for this module is MEM23004A Apply technical mathematics or equivalent. More details on this unit’s content from training.gov.au
The Definition of a 1 st and 2 nd Order Differential Equation
First order ordinary differential equations are of the formsee solution details here
Second order ordinary differential equations are of the form
where P, Q, R and G are continuous functions, see solution examples here.
Second order ordinary differential equations are analogous to the equations of the conic sections circle, ellipse, parabola and hyperbola obtained by the intersection of a plane with a cone. The general equation of the conic section has the form Ax 2 + Bxy + Cy 2 + Dx + Ey + F = 0, more details here.
Conic sections analogous to second order ordinary differential equations
Public domain image from the Wikimedia Commons
Second order ordinary differential equations also describe a wide range of real world phenomena such as electromagnetic, sound and fluid waves, heat transfer and forced vibrations.
Engineering Applications of 1 st and 2 nd Order Differential Equations
Electrical & Electronic Engineers use differential equations to model electromagnetic waves to determine interference effects from electromagnetic fields (EMF) and to design electrical machines such as electric motors.
Simulation of an electric motor used in electric vehicles and a wide range of industries to drive pumps, fans and machines, courtesy of Ansys Inc.
Mechanical Engineers use differential equations to model heat transfer in systems such as power stations and building heating and cooling.
Aerospace Engineers use differential equations to model the flow of air around aircraft.
Structural Engineers use differential equations to analyse the vibration of springs and dampers in structures, reducing the effects of earthquakes on buildings. Please watch the following videos.
Model of circular drum vibration.
Public domain image from the Wikimedia Commons
Chemical Engineers use differential equations to model chemical reactions and fluid flow in processes such as oil and metals refining.
Lessons on 1 st and 2 nd Order Differential Equations
To get a better understanding of differential equations review these lessons and quizzes from the following links:
Introduction to differential notation in physics - Interactive Mathematics
What is a differential equation? -Khan Academy
Solving Differential Equations by Separation of Variables- Interactive Mathematics
Separable differential equations - Khan Academy
Now go to the Interactive graphical tool from geogebratube.org
Related Rates problem- Interactive Mathematics
Look at Second order linear homogeneous differential equations- Interactive Mathematics
2 nd order linear homogenous differential equations - Khan Academy
All Khan Academy content is available for free at
Reference and Review Websites
Paul's Online Math Notes
MIT Open Courseware Differential equations
Interactive Maths exercises with solutions
SOS Mathematics exercises with solutions
Engineering Worked Examples
From Mathematics Resources for Education and Industry site, registration and login required here, to complete at least 2 of the exercises below.
From the integralmaths.org ‘My home’ home page:
Select
Mathematics Resources for Level 3 Engineering
Find a resource by mathematical content
Calculus
Then find the following lessons on;
RMS Values, Electrical/Electronic engineering
Student PDF (RMS values) complete this exercise
RMS values Interactive
Foam cushioning, Mechanical Engineering
Student PDF (Foam cushioning): complete this exercise
Three interactive resources provide a visual representation of impact discussed in the exercise.
Natural gas storage, Chemical Engineering
Student PDF (Natural gas storage): complete this exercise
On completion of the exercises you may wish to review the solutions contained in;
Teacher PDF (RMS values)
Teacher PDF (Foam cushioning)
Teacher PDF (Natural gas storage)
1 st and 2 nd Order Differential Equations Quiz
When you have completed the lessons, verify your knowledge with this multiple choice 1 st and 2 nd Order Differential Equation Quiz
If you need to repeat the quiz, click here (please note: you will not be able to repeat the first quiz)
PLEASE NOTE:
This Quiz is assessed and required as evidence that you have satisfactorily completed the requirements to receive RPL at your TAFE
Please attempt all questions on the quiz
To pass this quiz you will need to achieve at least 70%
You will have the option to email the results of the quiz to you, ensure that you save and print out these results as evidence that you have passed the module, your TAFE will ask to see these results.
An example of one of the many engineering applications of partial differential equations describes airflow, known as the Navier-Stokes equations. These equations are the basis for Computational Fluid dynamics or CFD used to model air flow around aircraft, ships and cars.
CFD airflow simulation around the 2009/2010 Basilisk Performance team's CO 2 dragster entry for F1 in Schools which won Best Engineered Car. CFD was also used by the Australian Cold Fusion team who won the world championship in 2012.
Original Image courtesy of Symscape licensed under the Creative Commons Attribution 3.0 Unported license.
Information
Definition
Application
Lessons
Example
Quiz
Course Information
Partial differential equations are part of the national engineering unit MEM23007A Apply calculus to engineering tasks. More details on unit content from training.gov.au
Required knowledge: Partial differential equations.
The prerequisite for this module are MEM23004A Apply technical mathematics or equivalent. More details on this unit’s content from training.gov.au
The Definition of a Differential Equation
A differential equation is partial if its solution depends on more than one variable and ordinary if its solution depends on a single variable.
For a function of two variable z = ƒ(x,y) we can keep one variable such as y fixed and treat our function ƒ as a function of x only, then calculate the derivative of f with respect to x (if it exists). Our new function is the partial derivative of ƒ with respect to x with the symbol .
Partial differential equations often model multidimensional systems such as systems in three dimensional spaces.
Engineering Applications
Electrical & Electronic Engineers use partial differential equations to model electromagnetic waves to determine interference effects from electromagnetic fields (EMF), and to design radio and microwave systems.
Simulation of microwave heating in an oven
Image from Adina R&D, Inc
Magnetic fields are used in the levitation of high speed trains, watch the video below on how high speed trains use magnetic fields.
They are also used in superconducting quantum levitation, watch this video on quantum levitation.
Mechanical Engineers use partial differential equations to model heat transfer in systems such as electronic products, vehicle brakes and building air conditioning.
Simulation of temperature and cooling air flow in an electronic product from Autodesk Simulation CFD.
Aerospace Engineers use partial differential equations to model the flow of air around aircraft prior to wind tunnel testing such as this Boeing 737 Max test, watch the video below.
Structural Engineers use partial differential equations to analyse the vibration of plates in structures, see the Resonance Experiment video below.
Structural vibration simulation from Autodesk
Lessons on Partial Differentiation
To get a better understanding of partial differential equations review these lessons:
MIT Open Courseware Introduction to Partial Differential Equations
Paul's Online Math Notes
Engineering Example
For a uniformly heat conducting cylindrical rod, the absolute temperature T at time t is given by T(x,t). Then T satisfies the partial differential equation
where,δ is a constant called the diffusivity of the rod material.
For the three dimensional case we would have a function with respect to the x, y, z coordinate system T(x,y,z,t) and T would satisfy the partial differential equation
In general, partial differential equations are much more difficult to solve analytically than ordinary differential equations. When an analytical solution is not possible numerical methods using computer software are often applied, see solution details here.
Partial Differentiation Quiz
When you have completed the lessons, verify your knowledge with this multiple choice Partial Differentiation Quiz
If you need to repeat the quiz, click here (please note: you will not be able to repeat the first quiz)
PLEASE NOTE:
This Quiz is assessed and required as evidence that you have satisfactorily completed the requirements to receive RPL at your TAFE
Please attempt all questions on the quiz
To pass this quiz you will need to achieve at least 70%
You will have the option to email the results of the quiz to you, ensure that you save and print out these results as evidence that you have passed the module, your TAFE will ask to see these results.
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13705 | https://pmc.ncbi.nlm.nih.gov/articles/PMC4067173/ | Eosinophil Granule Proteins: Form and Function - PMC
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Eosinophil Granule Proteins: Form and Function
K Ravi Acharya
K Ravi Acharya
From the ‡Department of Biology and Biochemistry, University of Bath, Claverton Down, Bath BA2 7AY, United Kingdom and
Find articles by K Ravi Acharya
‡,1, Steven J Ackerman
Steven J Ackerman
the §Department of Biochemistry and Molecular Genetics, College of Medicine, The University of Illinois, Chicago, Illinois 60607
Find articles by Steven J Ackerman
§,2
Author information
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From the ‡Department of Biology and Biochemistry, University of Bath, Claverton Down, Bath BA2 7AY, United Kingdom and
the §Department of Biochemistry and Molecular Genetics, College of Medicine, The University of Illinois, Chicago, Illinois 60607
1
To whom correspondence may be addressed. Tel.: 44-1225-386238; E-mail: K.R.Acharya@bath.ac.uk.
2
Recipient of a David Parkin Visiting Professorship from the University of Bath. To whom correspondence may be addressed: Dept. of Biochemistry and Molecular Genetics, the University of Illinois at Chicago, College of Medicine, 900 S. Ashland Ave., Chicago, IL 60607. Tel.: 312-996-6149; E-mail: sackerma@uic.edu.
Issue date 2014 Jun 20.
© 2014 by The American Society for Biochemistry and Molecular Biology, Inc.
PMC Copyright notice
PMCID: PMC4067173 PMID: 24802755
Abstract
Experimental and clinical data strongly support a role for the eosinophil in the pathogenesis of asthma, allergic and parasitic diseases, and hypereosinophilic syndromes, in addition to more recently identified immunomodulatory roles in shaping innate host defense, adaptive immunity, tissue repair/remodeling, and maintenance of normal tissue homeostasis. A seminal finding was the dependence of allergic airway inflammation on eosinophil-induced recruitment of Th2-polarized effector T-cells to the lung, providing a missing link between these innate immune effectors (eosinophils) and adaptive T-cell responses. Eosinophils come equipped with preformed enzymatic and nonenzymatic cationic proteins, stored in and selectively secreted from their large secondary (specific) granules. These proteins contribute to the functions of the eosinophil in airway inflammation, tissue damage, and remodeling in the asthmatic diathesis. Studies using eosinophil-deficient mouse models, including eosinophil-derived granule protein double knock-out mice (major basic protein-1/eosinophil peroxidase dual gene deletion) show that eosinophils are required for all major hallmarks of asthma pathophysiology: airway epithelial damage and hyperreactivity, and airway remodeling including smooth muscle hyperplasia and subepithelial fibrosis. Here we review key molecular aspects of these eosinophil-derived granule proteins in terms of structure-function relationships to advance understanding of their roles in eosinophil cell biology, molecular biology, and immunobiology in health and disease.
Keywords: Asthma, Crystal Structure, Enzyme Structure, Eosinophil, Inflammation, Eosinophil Granule Proteins
Introduction
Over the past three decades, the role of the eosinophil in human health and disease has received considerable attention (1, 2). The eosinophil has a vital role in allergic inflammatory processes that include asthma (3, 4). Evidence implicates the eosinophil and its granule proteins in host resistance to parasites, particularly helminths, but also antimicrobial activities toward bacterial, viral, and protozoan pathogens, and as mediators of hypersensitivity diseases. This evidence consists of associations between elevated levels of eosinophils in blood and the occurrence of disease; correlations between disease severity and degree of eosinophilia; findings that the cationic eosinophil-derived granule proteins (EDGPs)3 are toxic to cells and human tissues, producing changes mimicking those associated with disease (e.g. in bronchial asthma); deposition of the toxic EDGPs in diseased tissue; and observations that glucocorticoids suppress eosinophilia as part of their therapeutic effect (5).
The eosinophil is rich in cationic granule proteins, shows a striking respiratory burst with the production of toxic oxygen radicals, brominates tissue and proteins, presents antigen to T-cells, and expresses both hematopoietic and inflammatory cytokines (2). These findings support the eosinophil as a key player in allergic inflammation and tissue homeostasis in hypersensitivity diseases (1, 2). Eosinophil-deficient mouse models (3, 4, 6), including EDGP gene-deleted mice (7–9), provide unique insights into the role of the eosinophil, in both tissue damage/repair/remodeling and novel immunomodulatory roles that bridge host innate and adaptive immune responses in allergic and parasitic diseases.
Initial views of the eosinophil as providing direct antipathogen (helminth) activity are being supplanted by more nuanced views of the immunomodulatory functions of the eosinophil, including roles in supporting parasitic nematode survival (10, 11); Appleton and colleagues (10, 11) show that growth and survival of the muscle stage larvae of Trichinella spiralis require eosinophils, which promote accumulation of Th2 cells and inhibit the induction of inducible NOS by macrophages and neutrophils. Studies using a conditional eosinophil-deficient mouse strain (iPHIL) show that eosinophils modulate the immune and inflammatory character of inducible allergic responses in the lung (6). Clinical trials using anti-interleukin-5 (IL-5) antibodies to ablate eosinophils in the bone marrow and blood, as well as reduce tissue eosinophils in patients with the eosinophilic but not neutrophilic phenotype of asthma, show efficacy in reversing eosinophil-mediated tissue damage, remodeling, fibrosis, and airway dysfunction (12), as well as end-organ damage in the hypereosinophilic syndrome (13), highlighting the complex proinflammatory and immunomodulatory activities of the eosinophil in shaping the pathogenesis of these diseases.
Eosinophil-derived Granule Proteins
With expectations that understanding the properties, activities, and secretion of eosinophil proteins in disease states would provide insights into cellular function, the cationic components of the large crystalloid-containing specific (secondary) granule of the eosinophil (Fig. 1) were extensively studied. These proteins include major basic protein-1 and -2 (MBP-1, MBP-2), eosinophil peroxidase (EPX), eosinophil cationic protein (ECP), and eosinophil-derived neurotoxin (EDN) (Fig. 1). The Charcot-Leyden crystal protein/Galectin-10 (CLC/Gal-10), although not cationic, is a hydrophobic autocrystallizing protein comprising ∼7–10% of total eosinophil protein (Fig. 1). EPX and MBP-2 are the only EDGPs uniquely expressed by the eosinophil and not other cells; the other EDGPs are variably expressed in ∼10–100-fold lesser amounts by other blood leukocytes, tissues, and cells including basophils (MBP-1 (14, 15), EDN (16), CLC/Gal-10 (17, 18)), neutrophils (EDN, ECP) (19–21), liver (EDN) (22), and regulatory T cells (Tregs) (23, 24) and T H 2 central memory T cells (CLC/Gal-10) (25). Although these EDGPs are expressed and may be secreted by these other cell types at sites of host innate and adaptive immune responses and inflammation, aside from expression of CLC/Gal-10 by regulatory T cells (23), the roles and functions of these proteins beyond those of eosinophils have not been investigated, other than their potential use as biomarkers.
FIGURE 1.
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Structural representations of the human eosinophil granule proteins showing their location and known functions. The cationic granule proteins shown are: MBP-1 (Protein Data Bank (PDB) code 1H8U); EDN (RNase-2) (PDB code 1HI2); ECP (RNase-3) (PDB code 1QMT); and EPX/EPO) (based on molecular modeling using a myeloperoxidase structure, PDB code 1D2V). Also shown is CLC/Gal-10 (PDB code 1LCL), which is mainly cytosolic but also present in a small residual population of primary large core-less granules in the mature eosinophil. CLC/Gal-10 forms the hexagonal bipyramidal crystals (arrow) considered a hallmark of eosinophilic inflammation in tissues and body fluids in eosinophil-associated diseases.
Eosinophils form extracellular DNA traps (EETs), a component of innate antibacterial immune responses in a number of eosinophil-associated infectious, allergic, and autoimmune diseases (26, 27). EETs consist of a meshwork of DNA fibers (formed from mitochondrial rather than nuclear DNA), and a number of EDGPs, including MBP-1 and ECP, co-localize in the EETs, suggesting that they participate in trapping and killing of bacteria by this mechanism (28). EETs were identified in the lung in allergic asthma (29) and in a number of allergic/reactive skin diseases including allergic contact dermatitis, ectoparasitoses, and larva migrans (30). Multiple mechanisms of eosinophil priming and activation that include signaling through toll-like, cytokine, chemokine, and adhesion receptors initiate the formation of EETs containing the EDGPs (26, 31) in a process requiring activation of NADPH oxidase. Although the specific role of the EDGPs in EETs has not been determined, their formation may provide a mechanism for bringing together these cationic toxins with their pathogen targets, or alternatively limit collateral tissue damage by the EDGPs (26).
Eosinophil Major Basic Protein-1
MBP-1 (13.8 kDa, extremely basic, pI = 11.4) is an abundant granule protein localized to the electron-dense crystalloid core of the secondary granule. The protein is initially expressed as a 25.2-kDa polypeptide (pre-pro form), comprising a highly acidic “pro-domain” and highly basic MBP-1 protein. The pro-domain is thought to neutralize the basic nature of MBP-1 during its synthesis, transport, and packaging in the eosinophil granule, but structural or functional studies are lacking. MBP-1 is highly toxic to mammalian cells in vitro and can damage helminths, bacteria, and mammalian cells (32). MBP-1 damages cells by disrupting the lipid bilayer membrane (32) or altering the activity of enzymes within tissues. MBP-1 stimulates mediator (histamine) release from basophils and mast cells, activates neutrophils and platelets, and augments superoxide generation by alveolar macrophages. MBP-1 levels are elevated in biological fluids (e.g. sputa and bronchoalveolar lavage fluids) from patients with asthma and other eosinophil-associated diseases. MBP-1 induces bronchoconstriction when administered into monkey lung and induces transient airway hyperreactivity when administered into rat, rabbit, or monkey trachea (33). These effects could be neutralized by polyanions such as heparin or polyglutamic acids. MBP-1 has been implicated in tissue damage associated with eosinophil infiltration, notably to the respiratory epithelium in asthma. MBP-1 binds lactoferrin (34) and is a potent inhibitor of recombinant and eosinophil heparanase (35). Direct instillation of MBP-1 into the mouse airway induces airway remodeling, including induction of epithelial TGF-β and MMP1 expression and subepithelial fibrosis (36). In a model of Duchenne muscular dystrophy and the mdx mouse model of Duchenne muscular dystrophy, eosinophils were shown to lyse muscle cells in vitro through release of MBP-1, and MBP-1 promoted fibrosis of dystrophin-deficient muscle and attenuated the cellular immune response to these cells (37).
The role of MBP-1 in the development of allergen-induced pulmonary pathologies in asthma has been studied in mouse asthma models in which the MBP-1 gene was knocked out (MBP-1−/−) (38); MBP-1 deficiency had no effect on the development of allergen-induced airway histopathologies or inflammatory cell recruitment, nor any effects on airway hyperresponsiveness, which still developed in the absence of mouse MBP-1. Knock-out of the EPX gene (EPX−/−) and resulting EPX deficiency had no impact on asthma-associated pulmonary pathologies induced by allergen sensitization and provocation (39), suggesting that the contributions of MBP-1 and EPX to disease pathology in allergic diseases likely occurred via their combined activities. Crossing MBP-1−/− and EPX−/− knock-out mice to address this issue in double knockouts unexpectedly generated a novel strain of mice with a highly specific deficiency in eosinophilopoiesis, and therefore eosinophils (9). Although the mechanism for the eosinophil deficiency in these double knock-out mice remains to be determined, results suggest that: (i) granule protein gene expression and/or defective granulogenesis may be a checkpoint for survival of eosinophil progenitors, (ii) loss of concurrent expression of MBP-1 and EPX disrupts lineage-instructive gene regulatory mechanisms required for self-renewal or eosinophil progenitor survival, or (iii) most likely, dysfunctional granulogenesis leads to aberrant intracellular release of a toxicant such as mouse eosinophil-associated ribonucleases (EARs) capable of degrading intracellular RNAs, leading to the observed cell-autonomous defect (39).
A homologue of MBP-1, called MBP-2, was identified and shown expressed exclusively by eosinophils. MBP-2 has ∼66% amino acid sequence identity with MBP-1 but is significantly less basic (pI = 8.7) (40). In vitro activities of MBP-1 and MBP-2 appear similar, e.g. cell destruction, induction of superoxide anion, IL-8 release from neutrophils, and induction of histamine and leukotriene C4 release from basophils, but human MBP-1 is more potent than MBP-2 in these activities (52). MBP-2, present only in eosinophils, may be a useful biomarker for eosinophil-associated diseases, but its utility has not been evaluated.
MBP-1 is a monomer under physiologic conditions, but readily polymerizes in solution, forming insoluble aggregates due to the presence of five reactive thiol groups (in addition to four cysteines involved in disulfide bond formation) (41). MBP-1 is synthesized as a precursor that is proteolytically processed to the mature granule form during packaging into the crystalloid core of the granule. The pro-domain removed in this process is heavily glycosylated with N-glycans, O- glycans, and glycosaminoglycans, raising the molecular mass to ∼30–50 kDa. MBP-1 does not exhibit high sequence similarity to other proteins aside from weak similarity (23–28%) to C-type lectin domains and the low affinity IgE receptor FcϵRII (42).
The three-dimensional structure of MBP-1 (Fig. 1) (43) shows that its overall topology is similar to that of C-type lectin domains, in particular to lithostathine, a glycoprotein expressed by exocrine pancreas. We showed that none of the amino acid residues involved in calcium binding in classical C-type lectins is conserved in MBP-1. The region corresponding to the carbohydrate-binding site in MBP-1 is highly basic and thus differs in structure from that of the other C-type lectins. The crystal structure of MBP-1 in complex with heparin disaccharide (Fig. 2A) showed that heparan sulfate may be a ligand (in agreement with data showing MBP-1 binding to heparin) (44). It is likely that heparan sulfate is not the sole physiologic ligand for MBP-1, and it may have the capacity to recognize a wide variety of sulfated ligands.
FIGURE 2.
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Molecular details of the functional sites identified thus far.A, MBP-1 carbohydrate-binding region with bound heparin disaccharide and sulfate ions (PDB code 2BRS). B, EDN ribonucleolytic active site with bound inhibitor bis(adenosine)-5′-pentaphosphate (PDB code 1HI5). C, ECP ribonucleolytic active site with bound adenosine-2′,5′-diphosphate (PDB code 1H1H). D, CLC/Gal-10 carbohydrate recognition domain with bound mannose (PDB code 1QKQ).
MBP-1 has been shown to function in vitro and in vivo as an endogenous allosteric antagonist of the inhibitory muscarinic M2 receptor (45). Fryer and colleagues (45) showed that MBP-1 potently inhibits binding of N-methyl scopolamine (NMS) to guinea pig M2, but not M3, receptors. MBP-1 was found to inhibit atropine-induced dissociation of NMS-receptor complexes, showing that MBP-1 interaction with the M2 receptor is allosteric, suggesting that it may function as an endogenous allosteric inhibitor of agonist binding to this inhibitory receptor. Inhibition of NMS binding by MBP-1 was reversible by heparin, which binds and neutralizes MBP-1, consistent with structural findings. Because eosinophils secrete MBP-1 during allergen-induced airway inflammation, and treatment of allergen-challenged or ozone-challenged guinea pigs with heparin or a neutralizing antibody to MBP-1 restores M2 receptor function (46–48), eotaxin/CCR3-mediated eosinophil recruitment to and release of MBP-1 on airway nerves may contribute to M2 receptor dysfunction and vagally mediated bronchoconstriction in the asthmatic diathesis.
Eosinophil-derived Neurotoxin
EDN is a small, basic protein that belongs to the ribonuclease A (RNase A) superfamily (19). It is localized to the matrix of the secondary granule of the eosinophil (49) and is also known as RNase-2, nonsecretory RNase, and RNase-Us, the latter based on its specificity toward uridine-containing nucleotides. EDN is one of the most abundant RNases in humans and has been isolated from a wide variety of sources including eosinophil, placenta, liver, and urine. There is ∼3.3 μg of EDN/10 6 eosinophils (16). The initial identification of EDN was based on its induction of ataxia, incoordination, spasmodic paralysis, muscle stiffness, and killing of cerebellar Purkinje cells when injected intrathecally into rabbits, a paralytic syndrome termed the Gordon phenomenon (52, 53).
Although EDN shares 67% amino acid sequence identity with ECP, its sequence identity with RNase A is only 36%. The ribonucleolytic activity of EDN is ∼3–30-fold lower than that of RNase A, depending on the substrate (22). This enzymatic activity is a prerequisite for its cytotoxic, neurotoxic, and antiviral activities (22, 51, 54). The primary sequence of EDN contains a Trp-X-X-Trp motif between residues 7 and 10, specifying an unusual C-mannosylation of Trp-7. This involves attachment of an R-mannosyl residue via a C–C link to the indole moiety of Trp-7, the first example of this post-translational modification (55). Unlike the other EDGPs, EDN is a poor cationic toxin with limited toxicity for helminth parasites and mammalian cells at high concentrations (50, 56). However, as a ribonuclease, it is considerably more effective against single-stranded RNA viruses (51).
EDN activates human dendritic cells (DCs), leading to their expression of a variety of inflammatory chemokines, cytokines, growth factors, and soluble receptors (57). EDN also induces both phenotypic and functional maturation of DCs, as well as acts as an alarmin that activates the TLR2-MyD88 signaling pathway in DCs, enhancing Th2 immune responses (58).
The crystal structure of recombinant EDN (59) showed that the topology of the molecule includes the RNase A fold (Fig. 1) and that the core ribonucleolytic active site architecture (Fig. 2B) is conserved, although both ECP (60) and EDN exhibit significant differences at the peripheral substrate-binding sites (61). High-resolution crystal structures in complex with nucleotide inhibitors (59) present a detailed picture of differences and flexibility between EDN and RNase A in substrate recognition (Fig. 2B).
In the mouse, the related family of EARs, initially identified by Lee and colleagues (62), exhibits highly divergent properties both from one another and from human EDN and ECP. Zhang et al. (63) showed that there is a striking similarity between the evolutionary patterns of the mouse EAR genes and those of the major histocompatibility complex, immunoglobulin, and T cell receptor genes, enabling them to hypothesize that host defense and generation of diversity are among the primary physiological functions of the murine EARs. The discovery of a large number of divergent EARs suggests the intriguing possibility that these proteins have been specifically tailored through evolution to fight against distinct mouse pathogens (63).
Eosinophil Cationic Protein
ECP also belongs to the RNase A superfamily and is known as RNase-3 (20, 64, 65). Mature ECP is a small cationic polypeptide of 133 residues. Similar to EDN, it is located in the matrix of the specific granule of the eosinophil, but as compared with EDN (pI = 8.9), ECP is considerably more cationic (pI = 10.8). There is ∼5.3 μg of ECP/10 6 eosinophils (16). Like EDN, ECP induces the neurotoxic Gordon phenomenon (52). ECP has marked toxicity for a variety of helminth parasites, hemoflagellates, bacteria, single-stranded RNA viruses, and host tissues (66). Serum ECP levels can be used as a clinical tool for estimating eosinophil inflammatory activity in asthma and other allergic diseases, and levels are related to disease severity. The antibacterial activity and parasitic toxicity of ECP are greater than EDN (66). In vitro, ECP can function as an antiviral agent and may participate in host defense against the single-stranded RNA respiratory syncytial virus (67). The toxicity of ECP for bacteria and helminths does not appear related to its RNase activity (66), whereas RNase activity is required for its antiviral (67) and neurotoxic activities (52). The RNase activity of ECP is 100-fold lower than EDN for most RNA substrates, and their in vivo substrates have not been identified (20, 68).
The crystal structure of ECP (60) shows the “RNase fold” (Fig. 1), but with significant divergence from RNase A and EDN. The structure also shows how the cationic residues are distributed on the ECP surface, an observation that may have implications for understanding the considerable cytotoxicity of this enzyme. The structure of ECP in complex with adenosine-2′,5′-diphosphate revealed details of the active site (Fig. 2C) and a structural explanation for the lower substrate affinity and catalytic efficiency of ECP (69).
Although eosinophils and their EDGPs are associated with host defense responses against helminth parasites, a number of the EDGPs also possess antibacterial activity. ECP has antibacterial activities not shared by EDN (70, 71). Evidence is accumulating that eosinophils and the EDGPs may participate in host responses to certain bacterial infections (72). ECP is active in vitro against both Gram-negative and Gram-positive strains of bacteria, its mechanism of toxicity involving both the bacterial cell wall and the cytoplasmic membrane. Torrent et al. (73) propose and provide evidence for a novel molecular mechanism to explain the bacterial agglutinating activity of ECP, showing in situ formation of fibrillar, amyloid-like aggregates at the bacterial cell surface that bind amyloid diagnostic dyes. The agglutinating activity of ECP appears driven by the amyloid-like aggregation of the protein at the bacteria cell surface; elimination of the amyloidogenic behavior by a single point mutation (I13A) abolished both its agglutinating and its antimicrobial activities, the mutant being defective in triggering leakage and lipid vesicle aggregation. These findings support the novel concept that the amyloidogenic behavior of ECP, and possibly other EDGPs, participates in antibacterial host responses to infection, and suggest that the biophysical properties of bactericidal N-terminal peptides of ECP (amino acids 1–45) (74) or other EDGPs might guide development of novel antimicrobials (75).
Native ECP purified to homogeneity from blood leukocytes or purified eosinophils shows considerable molecular heterogeneity, from multiple glycosylated isoforms to the nonglycosylated native protein, as well as functional heterogeneity of these glycoforms relative to nonglycosylated ECP with respect to its cytotoxic activity for mammalian cells (76). A gene polymorphism in the coding region, ECP 434(G>C), determines the cytotoxicity of ECP for mammalian cells but has minor effects on fibroblast-mediated gel contraction (measure of fibroblast activation) and no effect on ECP RNase activity (77). This polymorphism changes an arginine (base at 434 is G) at position 97 to threonine (base at 434 is C). The ECP 434(G>C) polymorphism correlates with the natural course of Schistosoma mansoni infection (78) and with inflammatory bowel disease in an age- and gender-dependent manner (79). These ECP genotypes also show associations with the symptoms of allergy and asthma (80, 81). Venge and colleagues (82) reported that the various ECP glycoforms are processed during eosinophil secretion; the modifications to secreted ECP by activated eosinophils is explained in part by differences in their degree of N-linked glycosylation, such that secreted ECP acquires the masses of the more cytotoxic, less glycosylated, isoforms, including the nonglycosylated species (82), explaining in part the structural and functional heterogeneity of ECP as reported in the literature.
Eosinophil Peroxidase (EPX/EPO)
During activation, eosinophils can generate potentially toxic reactive oxygen species, which unlike the neutrophil, are mainly directed extracellularly (83). Oxidant production begins with the generation of superoxide by the membrane bound NADPH oxidase of eosinophils, which dismutates into hydrogen peroxide (H 2 O 2). EPX, the most abundant cationic protein of the matrix of the specific granule, uses this H 2 O 2 as an oxidizing substrate to generate potent oxidizing species, including hypohalous acids. In addition to bromide and chloride, EPX preferentially uses thiocyanate (SCN−) ions to generate HOSCN−, shown by Wang et al. (84, 85) to exert considerable biologic activity, e.g. as a potent oxidant inducer of tissue factor activity in endothelial cells; the HOSCN− generated by EPX from activated tissue eosinophils may induce the prothrombotic and proinflammatory endothelial and endocardial phenotypes responsible for thrombotic complications seen in the hypereosinophilic syndrome. EPX is structurally distinguished from the other EDGPs, being a two-chain hemoprotein (68 kDa); it is initially synthesized as a single chain precursor that is proteolytically processed to a 55-kDa heavy chain and 12.5-kDa light chain. EPX is highly cationic and similar to ECP and MBP-1 in this regard. Biochemical evidence suggests that EPX is structurally related to myeloperoxidase (Fig. 1) present in neutrophil-specific granules. Patients with myeloperoxidase deficiency have normal levels of EPX, indicating independent expression mechanisms for these peroxidases. There is ∼12 μg of EPX/10 6 eosinophils (16). EPX exerts some cytotoxic effects as a cationic toxin, being able to kill parasites (86, 87) and mammalian cells in the absence of H 2 O 2 and a halide co-factor. Furthermore, EPX exerts both anti-inflammatory (88) and pro-inflammatory (89) activities.
All of the human EDGPs including EPX itself, MBP-1, EDN, and ECP, are post-translationally modified by EPX via nitration at specific tyrosine residues during their synthesis and packaging in the developing eosinophil (90). High-resolution affinity-mass spectrometry showed single specific nitration sites at Tyr-349 in EPX and Tyr-33 in both EDN and ECP, and crystal structures of EDN and ECP, as well as structural models of EPX, suggest that these nitrated tyrosine residues are surface-exposed. Studies in EPX−/−, gp91phox−/−, and NOS−/− knock-out mice showed that tyrosine nitration of these cellular toxins and ribonucleases is mediated by EPX itself in the presence of H 2 O 2 and small amounts of nitrogen oxide. Thus, EPX appears to nitrate itself via an autocatalytic mechanism. Tyrosine nitration of the EDGPs was shown to occur during eosinophil differentiation and was independent of inflammation. The specific roles of EPX-mediated nitrosylation of the EDGPs in eosinophil-mediated innate host immune defense mechanisms characterized by their secretion during cell activation, e.g. against parasites or during tissue damage in parasitic infections or allergic responses, have not been determined.
Charcot-Leyden Crystal Protein
CLC protein forms distinctive bipyramidal hexagonal crystals, hallmarks of eosinophil participation in allergic and related inflammatory reactions. These crystals, found at sites of eosinophil infiltration in tissues and in body fluids and secretions, were identified more than 150 years ago. CLC is a small, slightly acidic (pI ∼5.1–5.7) 142-amino acid protein of 16.5 kDa (91). It was initially identified as an eosinophil lysophospholipase (LPLase) (92) but has since been assigned to the galectin superfamily of S-type animal lectins as the 10th member (Galectin-10) based on amino acid sequence (91), three-dimensional structure (CLC/Gal-10, Fig. 1) (93), and genomic organization (94). Importantly, we showed that CLC/Gal-10 lacks any LPLase activity and that the weak enzymatic activity initially associated with the purified protein was due to contamination by a highly active 75-kDa pancreatic LPLase also expressed by eosinophils (95).
The crystal structure of CLC/Gal-10 provides details of a carbohydrate recognition domain with both similarities to and important differences from other members of the galectin family (93). Structural studies showed that CLC/Gal-10 does not bind β-galactosides, but can bind mannose in the crystal in a unique manner different from the binding of lactosamine carbohydrates by other related galectins (Fig. 2D) (96). The physiologic significance of this mannose binding remains equivocal. The partial conservation of residues involved in carbohydrate binding leads to significant changes in the topology and chemical nature of the carbohydrate recognition domain and has implications for glycan recognition by CLC/Gal-10 in vivo. These findings are beginning to provide clues toward identifying a physiologically relevant ligand for CLC/Gal-10 and understanding its potential intracellular and extracellular roles in eosinophil biology.
Our recent experiments using Southwestern (ligand) blotting, co-purification, co-immunoprecipitation, and confocal microscopy show that CLC/Gal-10 interacts in vitro and intracellularly in activated eosinophils with the two glycosylated human eosinophil granule cationic ribonucleases, EDN and ECP.4 Studies in human blood eosinophils using immunofluorescence confocal microscopy show that interferon-γ activation induces the rapid co-localization of CLC/Gal-10 with EDN and CD63.4 Because CLC/Gal-10 does not inhibit the ribonuclease activity of EDN, it may function instead as a carrier for the sequestration and vesicular transport of these ribonucleases, during granulogenesis in the differentiating eosinophil, and during piecemeal degranulation (Fig. 3) in the activated eosinophil, enabling their extracellular functions in host defense and allergic inflammation without intracellular damage to the eosinophil itself during their secretion.
FIGURE 3.
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Mechanisms of eosinophil degranulation. Eosinophils release their secondary (specific) granule contents, including the granule cationic proteins and other immunomodulatory mediators, by four different mechanisms. In classical exocytosis, mainly seen at sites of bacterial infection, the contents of single granules are released by fusion of the granule membrane with the plasma membrane lipid bilayer. In compound exocytosis, mainly seen with eosinophil degranulation onto helminth parasites, a number of granules first coalesce and fuse, and the cationic protein contents are then released through a single fusion pore at the plasma membrane. In piecemeal degranulation, mainly seen in eosinophil inflammatory responses in human tissues, secretory vesicles bud from granules, capturing granule matrix and/or core contents, and are targeted to the plasma membrane in a fashion analogous to the release of neurotransmitters from neurons. Differential secretion of the cationic proteins from the matrix (EPX, EDN/RNase-2, ECP/RNase-3) or core (MBP-1, MBP-2) of the granule has been shown to occur, leaving coreless granules with intact matrix or intact cores with no matrix in the cell. In cytolysis, intact whole granules are deposited in tissues and body fluids after disruption of the plasma membrane due to eosinophil necrosis.
In addition to its considerable expression at both mRNA and protein levels in eosinophils and basophils, CLC/Gal-10 was identified as a constituent of human regulatory T cells. A global proteomics analysis of highly purified human CD4+CD25+ Tregs identified CLC/Gal-10 as a novel biomarker, shown by siRNA knockdown to be essential for maintaining Treg anergy and suppressive functions on T cell activation (23). The mechanism by which CLC/Gal-10 participates, through galectin-type interaction with a glycan ligand(s) or protein-protein interaction, in maintaining the CD4+CD25+ Treg phenotype has not been established.
A number of studies identify CLC/Gal-10 as a potentially useful biomarker of eosinophil involvement in asthma, allergic rhinitis, and other eosinophil-associated diseases. Elevated levels of CLC/Gal-10 have been measured by one of us, by ELISA, as a biomarker of active eosinophilic inflammation that is highly correlated with the number of tissue (esophageal) eosinophils in eosinophilic esophagitis (97), a rare immune-mediated food allergic disease of increasing incidence. CLC/Gal-10 was also identified as a potentially useful biomarker of eosinophilic airway inflammation in induced sputum for identifying the eosinophilic phenotype of asthma to guide treatment considerations (98). In celiac disease, CLC/Gal-10 expression was found related to both disease activity (histologic grade) and numbers of tissue eosinophils in intestinal lesions, suggesting it as a novel biomarker for evaluating tissue damage and eosinophil involvement in the pathogenesis of this gluten intolerance. Finally, genetic variations (SNPs) in the promoter region of the CLC gene were identified as potential susceptibility biomarkers for allergic rhinitis (99), with the pattern of variation compatible with a recessive inheritance model and observed increased levels of CLC/Gal-10 protein in the nasal fluid of patients with allergic rhinitis during the allergy season (100).
Perspectives
With a goal toward understanding the roles and specific functions of the EDGPs in the normal and pathologic activities of the eosinophil, significant progress has been made by determining the three-dimensional structures of four of these mediators using x-ray crystallographic approaches. These studies provided novel insights and vital clues toward understanding the structural basis for their unique biologic and enzymatic activities at a molecular level, paving the way for future “form and function” analyses to better define their unique biochemical properties, biologic activities, and pathologic contributions to eosinophil-mediated inflammatory responses, tissue damage and repair, remodeling, and fibrosis, as well as innate and adaptive host immune responses to infectious agents including parasitic helminths, bacteria, and viruses.
Acknowledgments
We thank Drs. Nethaji Thiyagarajan and Geoffrey Masuyer for Figs. 1 and 2.
4
C. B. Doyle and S. J. Ackerman, unpublished results.
3
The abbreviations used are: EDGP
eosinophil-derived granule protein
MBP
major basic protein
EPX
eosinophil peroxidase
ECP
eosinophil cationic protein
EDN
eosinophil-derived neurotoxin
CLC/Gal-10
Charcot-Leyden crystal protein/Galectin-10
Treg
regulatory T cell
EETs
eosinophil extracellular DNA traps
EAR
eosinophil-associated ribonuclease
NMS
N-methyl scopolamine
DC
dendritic cell
LPLase
lysophospholipase.
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13706 | https://www.cuemath.com/calculators/n-choose-k-calculator/ | n Choose k Calculator
A combination is a mathematical technique that determines the number of possible arrangements in a collection of items where the order of the selection does not matter.
What is n Choose k Calculator?
'n Choose k Calculator' is an online tool that assists in calculating the number of possible combinations of selecting a sample of k elements from a set of n distinct objects. Online n Choose k Calculator helps you to calculate the number of combinations in a few seconds.
n Choose k Calculator
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Please follow the steps below on how to use the calculator:
How to Find the Combinations?
The combinations are defined as the number of ways in which a sample of r elements can be selected from n distinct objects that's why it is also referred to as 'n choose k'. To determine the number of combinations the following formula is used:
C(n, k) = n!/(k!(n - k)!)
It is read as the number of possible combinations of selecting a sample 'k' from 'n' distinct objects.
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Example 1:
Find the number of ways in which 6 balls can be selected from a bag containing 9 different colored balls
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Required Sample size = 6
Number of combinations = C(n,k) = n!/(k!(n-k)!)
C(9,6) = 9!/(6!(9-6)!)
C(9,6) = 9!/(6!(3)!)
C(9,6) = 84.
Therefore, the total number of combinations to select 6 balls from a bag of 9 balls is 84.
Example 2:
Find the number of ways in which 3 balls can be selected from a bag containing 7 different colored balls
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Total number of balls = 7
Required Sample size = 3
Number of combinations = C(n, k) = n!/(k!(n - k)!)
C(7,3) = 7!/(3!(7-3)!)
C(7,3) = 7!/(3!(4)!)
C(7,3) = 35.
Therefore, the total number of combinations to select 3 balls from a bag of 7 balls is 35.
Example 3:
Find the number of ways in which 10 balls can be selected from a bag containing 15 different colored balls
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Total number of balls = 15
Required Sample size = 10
Number of combinations = C(n,k) = n!/(k!(n-k)!)
C(15, 10) = 15!/(10!(15-10)!)
C(15, 10) = 15!/(10!(5)!)
C(15, 10) = 3003
Therefore, the total number of combinations to select 6 balls from a bag of 9 balls is 84.
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13707 | https://scholarworks.umt.edu/cgi/viewcontent.cgi?article=1526&context=tme | Published Time: Thu, 16 Dec 2021 18:30:38 GMT
The Mathematics Enthusiast The Mathematics Enthusiast
Volume 18
Number 1
Numbers 1 & 2 Article 19
1-2021
The Suuji Approach to Multi-Digit Addition: Using Length to The Suuji Approach to Multi-Digit Addition: Using Length to
Deepen Students’ Understanding of the Base 10 Number System Deepen Students’ Understanding of the Base 10 Number System
Ryota Matsuura
Olaf Hall-Holt
Nancy Dennis
Michelle Martin
Sarah Sword
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Recommended Citation Recommended Citation
Matsuura, Ryota; Hall-Holt, Olaf; Dennis, Nancy; Martin, Michelle; and Sword, Sarah (2021) "The Suuji
Approach to Multi-Digit Addition: Using Length to Deepen Students’ Understanding of the Base 10
Number System,"
The Mathematics Enthusiast : Vol. 18 : No. 1 , Article 19.
DOI:
Available at:
This Article is brought to you for free and open access by ScholarWorks at University of Montana. It has been
accepted for inclusion in The Mathematics Enthusiast by an authorized editor of ScholarWorks at University of
Montana. For more information, please contact scholarworks@mso.umt.edu .TME, vol. 18, nos.1&2, p.309
The Mathematics Enthusiast, ISSN 1551-3440, vol. 18, nos.1&2 , pp.309-324 2021© The Author(s) & Dept. of Mathematical Sciences-The University of Montana
The Suuji Approach to Multi-Digit Addition:
Using Length to Deepen Students’ Understanding of the Base 10 Number System
Ryota Matsuura 1
St. Olaf College
Olaf Hall-Holt
St. Olaf College
Nancy Dennis
Prairie Creek Community School
Michelle Martin
Prairie Creek Community School
Sarah Sword
Education Development Center, Inc.
Abstract : We describe the Suuji representation of numbers which aims to deepen elementary students’ understanding of the base 10 system. (“Suuji” means “number” in Japanese.) This representation takes a two pronged approach of (1) making the place value more explicit and (2) using length to represent numbers, thus allowing students to reason spatially. We taught multi-digit addition using the Suuji representation to 20 second and third grade students. The article uses lesson descriptions and student work to illustrate the Suuji approach, as well as its impact on student learning.
Keywords : Base 10 number system, Multi-digit addition, Place value, Measurement-based approach, Spatial reasoning, Representation, Physical manipulatives.
Introduction
Students typically learn multi-digit addition in second or third grade (CCSSI, 2010). Some struggle with this work when they do not have deep understanding of place value in the base 10 number system. For instance, a student might consider 37 as thirty-seven ones, rather than as three tens and seven ones. And when faced with 37 + 56, a student might resort to counting by ones ( figure 1 ).
1
matsuura@stolaf.edu Matsuura et al. p.310
Figure 1: Student adds by counting by ones.
To alleviate this and other related issues, we developed the Suuji representation that makes each place value explicit. (“Suuji” means “number” in Japanese.) For example, H3 R1 S4, which is a short-hand for 3 houses, 1 room, and 4 shelves, corresponds to the base 10 number 3.14.
We taught multi-digit addition using the Suuji representation to 20 second and third grade students at a rural public school. In 2016–2017, the school had 180 students, with 14% minority, 16% free/reduced lunch, and 17% special education students. These students had a wide range of prior experiences. Some added by counting, without considering place values ( figure 1 ). Others proficiently added multi-digit whole numbers.
The article uses lesson descriptions and student work to illustrate the Suuji approach. We share pre-and post-assessment results to describe effects on student learning. While we used decimal numbers with our students, the article also describes how our approach may be adjusted for whole number addition.
Note: We developed physical manipulatives, partner games, and exercises through which students engaged with the Suuji representation. For sample lesson plans, as well as templates for these manipulatives, contact the first author via email.
Background
The Suuji representation aims to deepen students’ understanding of numbers by making the place value more explicit. In this regard, our approach resembles often-used manipulatives such as the Base Ten blocks and the abacus (e.g., Cotter, 2000). However, the Suuji approach distinguishes itself by TME, vol. 18, nos.1&2, p.311
representing numbers as lengths , allowing students to reason spatially. More specifically, the Suuji blocks (described in Lesson 1 ) emphasize height as the underlying physical quantity of interest.
The Suuji representation was inspired, in part, by the work of Russian psychologist Vasily Davydov. During the 1960s, Davydov and his colleagues developed an early elementary mathematics curriculum based on the measurement notion of numbers, rather than the traditional approach that uses counting (Davydov, 1990). In Davydov’s curriculum, children study scalar quantities such as length, area, and volume, which can be experienced visually and tactilely (Schmittau, 2005). For example, early first grade students might compare two lengths and make them equal by adding to the smaller or subtracting from the larger. Shown below is a schematic diagram of their thinking (Schmittau, 2005, p. 19):
Figure 2: Measurement notion of numbers.
Davydov’s approach provides opportunities for early algebraic thinking (Bass, 2015). Various studies, including those conducted in American schools, have shown that his approach is effective in fostering young children’s mathematical understanding (Schmittau, 2003; Venenciano, Slovin, & Zenigami, 2015). Suuji mathematics, with its emphasis on length, espouses this measurement-based approach, bringing coherence to the mathematics that students learn and providing them access to algebraic thinking at an early age.
The Suuji representation takes the “best of both worlds” route by combining Davydov’s measurement-based approach with an explicit emphasis on place value. As described in Lesson 2 , the Suuji blocks allow exchanges between ones and tens, or tens and hundreds; but in distinction from the Base Ten blocks, these exchanges occur in the context of towers of equal height. Thus, students can use Matsuura et al. p.312
not only their grasp of discrete exchanges, but also their spatial reasoning with lengths, to undergird their learning about the base 10 number system.
The Suuji representation was developed in collaboration with classroom teachers, who were looking for more than what their students were getting out of the Base Ten blocks. The teachers chose the representation using neighborhoods, houses, rooms, and shelves, because they thought it would be fun and relevant to their students. The key here is not “houses” or “rooms,” however. Instead, it is the way in which the Suuji representation allows students to visualize and work with numbers and number relationships in a meaningful way.
In teaching these lessons, our goal was not simply to have students perform multi-digit arithmetic. Rather, we designed activities that allow them to develop richer understanding of the base 10 system
through multi-digit addition, even as they gain fluency and understanding of multi-digit operations. Using the Suuji blocks, students learn to name numbers (such as H3 R1 S4 that we saw in the introduction), compare numbers, rewrite numbers (analogous to grouping in base 10), and add numbers. We view and treat these concepts not as separate skills to acquire, but all part of an interconnected whole. In this regard, our approach and this paper serve as an illustration of Skip Fennell’s rendering of how counting, place value, comparing and ordering, and operations are interconnected (Fennell, 2015).
Lesson 1: Towers and Tower Names
We describe five lessons on Suuji addition, each lasting 40 minutes. Lesson descriptions have been edited for clarity.
In this first lesson, students learned about towers and tower names . Towers are made from building blocks that correspond to powers of 10. The block types and corresponding decimals are shown:
● Neighborhoods tens
● Houses ones
● Rooms tenths
● Shelves hundredths TME, vol. 18, nos.1&2, p.313
For example, a tower containing 2 houses (2 ones) and 3 rooms (3 tenths) corresponds to 2.3. We did not reveal this tower-decimal correspondence until later.
Note: This correspondence relies on general place value structure, rather than on particular place values. For example, if we had been working with whole numbers instead of decimals, we would have had neighborhoods, houses, rooms, and shelves correspond to thousands, hundreds, tens, and ones, respectively.
We began the lesson by showing the two towers in figure 3 . Blocks are stacked in right-alignment. Tower name describes how its blocks are arranged from the bottom. For example, the tower R1 H1 R2 H1 (figure 3 , right) has one room at the bottom, followed by one house, two rooms, and one house. In this lesson, students worked with towers containing houses and rooms only.
Students played a partner game. Each student received a set of houses and rooms made of foam boards and cards with tower names. During each round, students built the tower whose name was written on the card. Then they studied their partner’s tower and wrote its name. This gave students practice in converting between towers and tower names. After each round, partners compared their towers to determine which is taller. Though we did not tell the students, the cards were designed so that partners built towers of equal height in each round (e.g., H2 R3 and R1 H1 R2 H1).
Figure 3: Towers H2 R3 and R1 H1 R2 H1.
Afterwards, we discussed whether or not towers with different names could have equal height .Students offered these ways to prove that two towers have equal height: Matsuura et al. p.314
● The two towers contain the same number of blocks for each type. Example: H2 R3 and R1 H1 R2 H1 both contain two houses and three rooms, so their heights are equal.
● We can rearrange the blocks in one tower to form the other. Example: Students said R1 H1 R2 H1 is “wobbly.” Moving its houses to the bottom forms the more stable H2 R3.
During the game, partners built towers H1 and R10. Students observed that 1 house and 10 rooms have equal height, and we wrote H1 = R10. In the Suuji unit, the equal sign indicates equality of heights. The towers themselves need not be the same. Although we did not delve further into this issue with our students, this way of using the equals sign sets up a reasonable way of relating two objects as equivalent but not necessarily “the same.”
Lesson 2: Standard Towers
This lesson focused on standard towers. A tower is standard if (1) its larger blocks are below the smaller blocks and (2) it uses as few blocks as possible. To standardize H2 R8 H1 R5, for example, rearrange its blocks with the largest type (houses) at the bottom. This yields H3 R13. Then reduce the number of blocks in this tower, because 13 rooms have height equal to one house and three rooms (R13 = H1 R3). Combining these yields H4 R3 ( figure 4 ). Standardizing does not change a tower’s height, since the only operations involved are rearranging blocks and exchanging 10 blocks with one block of equal height (i.e., 10 rooms with 1 house). TME, vol. 18, nos.1&2, p.315
Figure 4: Towers H2 R8 H1 R5, H3 R13, and H4 R3.
We began the lesson by recalling the pair of towers with equal height, H2 R3 and R1 H1 R2 H1 (figure 3 ). This lesson’s theme was an “invariant,” or a property that does not change. Students would manipulate a tower to look like another without changing its height—thus, the height is an invariant.
Students played another partner game. In each round, the “builder” built two towers of equal height whose names were written on a card. Then the “manipulator” manipulated the first tower (tower A) to look like the second tower (tower B). They switched roles after each round. By design, tower B was always standard—so, students standardized tower A in each exercise. Standardizing involved rearranging houses and rooms, and exchanging 10 rooms with 1 house. The rounds were sequenced so that students first grappled with individual skills of rearranging and exchanging, then gradually combined them.
During the whole class discussion afterwards, we created H2 R8 H1 R5 and issued a challenge: “Make a tower of equal height that uses as few blocks as possible.” Students suggested rearranging the houses to the bottom, then exchanging 10 rooms with 1 house, obtaining H4 R3. Then, we formally Matsuura et al. p.316
introduced the notion of standard towers and pointed out that students had made standard towers throughout the partner game.
Although we did not yet connect towers and decimals, some students noticed that houses and rooms behaved like base 10 numbers.
Lesson 3: Standardizing Tower Names
Students began the transition from visual to symbolic, to standardize towers using only tower names, without the visual support of blocks. We also introduced a new block type called shelves .
We began by showing students the visual image in figure 5a . We said, “This middle tower is taller than H1 R5 but shorter than H1 R6. The four tiny blocks are shelves .” We discussed how 1 room and 10 shelves have equal height (R1 = S10), just like 1 house and 10 rooms (H1 = R10).
Figure 5a (left): Towers H1 R5, H1 R5 S4, and H1 R6.
Figure 5b (right): Zoomed-in look at the top of the tower H1 R5 S4.
We demonstrated how to standardize H1 R5 S4 H1 R8 S7 using tower names only ( figure 6 ). Students suggested combining houses, rooms, and shelves, with the largest blocks at bottom—this yielded H2 R13 S11. We called such a tower a stable tower and explained that it is a helpful step before standardizing. Then we exchanged the 11 shelves with 1 room and 1 shelf, and the 14 rooms (13 original plus 1 extra) with 1 house and 4 rooms, obtaining H3 R4 S1. TME, vol. 18, nos.1&2, p.317
Figure 6: Standardizing H1 R5 S4 H1 R8 S7 using tower names.
Students worked on a series of exercises in which they saw a tower name like H1 R5 S4 H1 R8 S7 and found the name of the standardized tower, i.e., H3 R4 S1. The exercises were scaffolded by gradually increasing the complexity of towers. In the early exercises, students were also shown images of tower blocks with tower names. When students were ready, we removed the visual support, and students worked purely symbolically. Many students devised their own approaches for standardizing ( figures 7a and 7b ).
Figure 7a: Standardizing H2 R17. Figure 7b: Standardizing H1 R8 H1 R3. Matsuura et al. p.318
Lesson 4: Tower Addition
Students learned to add towers. Example: To find H2 R8 + H1 R5, stack the two towers on top of each other, obtaining a new tower H2 R8 H1 R5; standardizing that yields the sum H4 R3 ( figure 8 ).
Figure 8: Tower sum H2 R8 + H1 R5 = H4 R3.
Since students had previously standardized towers such as H1 R5 S4 H1 R8 S7 (which corresponds to the sum H1 R5 S4 + H1 R8 S7), tower addition felt familiar. This lesson also provided more experience with standardizing.
Students worked on a series of exercises in which they saw two tower names and found the sum. Again, the exercises were scaffolded to introduce sums of gradually increasing complexity. Tower blocks were shown in early exercises, and this visual support was later removed.
During the whole class debriefing, we used the sum H8 R2 + H6 R1, where the two towers together contained more than 10 houses, to introduce a new block type called neighborhoods . 1 neighborhood and 10 houses have equal height (N1 = H10), thus H8 R2 + H6 R1 = H14 R3 = N1 H4 R3.
Note: Tower addition is essentially regrouping. For example, to find H2 R8 + H1 R5, first create a stable tower H3 R13; this is analogous to combining the digits in each place. Then exchange the 13 rooms with 1 house and 3 rooms to obtain H4 R3; this exchange corresponds to the regrouping process. TME, vol. 18, nos.1&2, p.319
Lesson 5: Big Reveal
We introduced a shorthand for tower names called decimal numbers . Example: Decimal number 2.5 corresponds to tower name H2 R5. Each digit in the decimal number corresponds to a block type, and the dot separates houses and rooms. Other examples include 1.54 = H1 R5 S4 and 1.87 = H1 R8 S7. Thus, a decimal sum 1.54 + 1.87 corresponds to H1 R5 S4 + H1 R8 S7. This equals H3 R4 S1 (see figure 6 ), whose decimal form is 3.41. Hence, 1.54 + 1.87 = 3.41.
To summarize, computing a decimal sum like 1.54 + 1.87 involves:
Converting the decimals into tower names. Here, 1.54 + 1.87 becomes H1 R5 S4 + H1 R8 S7.
Computing the tower sum, i.e., H1 R5 S4 + H1 R8 S7 = H3 R4 S1. The sum must be standard.
Converting the standardized tower sum back to decimal form, i.e., H3 R4 S1 becomes 3.41.
Students worked on a series of exercises in which they saw a pair of decimals and computed the sum. Initially, the exercises guided students through the three steps above. Eventually, this support was removed—students were simply given, say, 2.78 + 5.63, and computed the sum ( figure 9 ).
Figure 9: Computing 2.78 + 5.63.
Effects on Student Learning
To measure effects on student learning, we administered pre- and post-assessments. They were administered three weeks apart, at the beginning and end of a larger unit that contained the five lessons. The assessments contained the same ten addition problems such as 37 + 56 and 5.76 + 3.37. On the post-assessment, students demonstrated their understanding of the role of place value in addition. Figure 10
shows the work of the student who, on the pre-assessment, computed 37 + 56 by counting ( figure 1 ). Matsuura et al. p.320
Here, he thinks about the place of each digit in 37 and 56 by grouping them as neighborhoods and houses. This conceptualization is very different from where he started.
Figure 10: Computing 37 + 56 on the post-assessment.
On the post-test, 84% of students (16 out of 19) comfortably added three- and four-digit numbers, as opposed to 10% (2 out of 20) on the pre-test. We saw fluidity in their thinking—students readily transferred their understanding of two-digit addition to the four-digit case. If students are merely building procedural knowledge, four-digit numbers pose more challenge than two-digits, because procedures become more complex. In Suuji addition, more digits did not increase difficulty for students, suggesting that they built conceptual understanding of multi-digit addition.
We gave students opportunities to develop their own understanding, which is empowering. Several students who had previously lacked confidence in mathematics said they felt they are now good at math. A student who usually struggled mathematically said, “I like this Suuji math. It’s just at my level.” He was doing far more complex addition than he had been doing prior to the Suuji unit.
Suuji Subtraction
After the Suuji unit, we worked with one of the students on subtraction. We asked her to compute 823 – 286, although she had not subtracted three-digit numbers before. She tried converting 823 to a tower name, but realized that another block type was needed to represent the 8 in 823. She chose
boroughs , where 1 borough and 10 neighborhoods have equal height (B1 = N10). Thus, 823 – 286 became B8 N2 H3 – B2 N8 H6. She could not subtract H6 from H3, however, so she exchanged one of TME, vol. 18, nos.1&2, p.321
the neighborhoods in B8 N2 H3 with 10 houses; then she exchanged one of its boroughs with 10 neighborhoods. Thus, B8 N2 H3 became B7 N11 H13, after which she subtracted ( figure 11 ).
Figure 11: Computing 823 – 286.
As this episode illustrates, Suuji lessons provided students with a sense of ownership and empowerment—this particular student applied her understanding of Suuji concepts to solve an unfamiliar problem. We explicitly discussed with all students that the Suuji approach was designed to help them understand how numbers worked, and this gave them license to think about extending the model. If we had presented it as a fait accompli , they may not have felt as empowered to tinker and explore.
Concluding Remarks: Benefits of the Suuji Approach
In this section, we describe the various benefits of the Suuji approach in deepening students’ understanding of numbers and the base 10 representation.
The Suuji representation such as H3 R1 S4 makes each place value explicit. Cotter (2000) notes how the English language—with words such as eleven , twelve , and thirteen —blurs the patterns for counting and obscures the groupings that occur in base 10. In contrast, many Asian languages have predictable counting patterns. For example, the number 37 is named 3-ten 7 in Japanese (san-ju shichi, where san, ju, and shichi mean three, ten, and seven, respectively). The Suuji representation follows Cotter’s recommendation of learning the base 10 system with an approach that highlights such counting patterns.
We have discussed (in the Background section) how Suuji mathematics takes the “best of both worlds” route by combining Davydov’s measurement-based approach with an explicit emphasis on place Matsuura et al. p.322
value. This connection is highlighted by the fact that there is a corresponding visual depiction of numbers, in the form of towers. Thus, students can see concretely that, for example, H3 R13 and H4 R3 have equal height, which offers a visual confirmation for the grouping that occurs. Such a visual cue is a powerful way to experience the relationship between these number representations.
Games and exercises played a prominent role in the Suuji unit. They provided targeted practice with individual skills, while gradually increasing the complexity of the tasks; this also allowed students to work at their own pace. Students worked with many examples, which gave them opportunities to “look for and express regularity in repeated reasoning” (CCSSI, 2010, p. 8). As a result, our students found general methods, such as their own approaches to standardizing.
An essential aspect of our approach is the set of physical manipulatives that allowed students to see and interact with tower blocks. The tactile nature of physical towers made these houses, rooms, and shelves feel concrete and real, even when working with their symbolic counterparts. Moreover, students connected with the playful nature of the Suuji representation. Learning about the base 10 system can be overwhelming and dry for many children. But our students enjoyed building towers with houses and rooms, seeing them get “wobbly,” and rearranging and exchanging blocks to make the towers standard. And it is precisely this “playing” with the tower blocks that led students to figure out for themselves and make sense of grouping in base 10.
Our students worked with decimal numbers, although (as described in Lesson 1 ) the Suuji representation can be adjusted to work with whole numbers only. Decimal numbers intimidate many students. They think they have to relearn everything about integers—like how to add them—in decimal setting. By treating decimals as a separate topic, we balkanize students’ mathematical understanding. We are telling students there are “special cases” they will learn later. We developed the Suuji approach with a belief that learning about decimals and multi-digit addition together brings coherence to students’ understanding. We can teach fluidity with the base 10 number system. When learning to add multi-digit numbers, rather than working exclusively with integers, students can learn to view integers (and decimals) as part of a continuum in their understanding about numbers. The Suuji approach gets at the core of how TME, vol. 18, nos.1&2, p.323
addition is done in the base 10 system, without separating the process into different cases according to different types of numbers.
As students engaged with Suuji mathematics, they bumped into very significant mathematical ideas, used in virtually all areas of mathematics: equivalence, invariance, underlying structures of the base 10 number system, multiple representations and going back and forth between them, and the value of “unpacked” representations and “compact” (or standardized) representations for working with and communicating about mathematical ideas. It is convenient to build a tower by simply stacking two towers on top of each other (e.g., H2 R8 H1 R5 in figure 8 ). But when communicating about numbers, it is convenient to have a standardized, compact way to describe them (e.g., H4 R3 in figure 8 ).
Lastly, we did not present the Suuji representation as a finished product. Instead, we gave students free rein to extend the model and use it to ask questions beyond the scope of what we covered in class—in essence, using the model as an authentic mathematical tool, thus allowing students to be “tinkerers” and “inventors” (Cuoco, Goldenberg, & Mark, 1996).
References
Bass, H. (2015). Quantities, numbers, number names, and the real number line. In X. Sun, B. Kaur, & J. Novotna (Eds.), Proceedings of the twenty-third ICMI Study: Primary Mathematics Study on Whole Numbers ( pp. 10–20). Macau, China: University of Macau.
Common Core State Standards Initiative [CCSSI]. (2010). Common Core Standards for Mathematics. Washington, DC: National Governors Association Center for Best Practices and the Council of Chief State School Officers.
Cotter, J. A. (2000). Using language and visualization to teach place value. Teaching Children Mathematics, 7 (2), 108–114. Matsuura et al. p.324
Cuoco, A., Goldenberg, E. P., & Mark, J. (1996). Habits of mind: An organizing principle for mathematics curriculum. Journal of Mathematical Behavior, 15 (4), 375–402.
Davydov, V. V. (1990). Types of generalisation in instruction: Logical and psychological problems in the structuring of school curricula . Reston, VA: NCTM. (Original published in 1972)
Fennell, F. (2015, January 5). Critical Foundations . Retrieved from
Schmittau, J. (2003). Beyond constructivism and back to basics: A cultural historical alternative to the teaching of the base ten positional system. In B. Rainforth & J. Kugelmass (Eds.), Curriculum and instruction for all learners: Blending systematic and constructivist approaches in inclusive elementary schools (pp. 113–132). Baltimore, MD: Brookes Publishing.
Schmittau, J. (2005). The development of algebraic thinking: A Vygotskian perspective. ZDM , 37 (1), 16– 22.
Venenciano, L., Slovin, H., & Zenigami, F. (2015, June). Learning Place Value through a Measurement Context . Paper presented at the twenty-third ICMI Study: Primary Mathematics Study on Whole Numbers, Macau, China. |
13708 | https://pmc.ncbi.nlm.nih.gov/articles/PMC5902026/ | Treatment of late sequelae after radiotherapy for head and neck cancer - PMC
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Published in final edited form as: Cancer Treat Rev. 2017 Jul 18;59:79–92. doi: 10.1016/j.ctrv.2017.07.003
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Treatment of late sequelae after radiotherapy for head and neck cancer
Primož Strojan
Primož Strojan
aDepartment of Radiation Oncology, Institute of Oncology, Ljubljana, Slovenia
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a,, Katherine A Hutcheson
Katherine A Hutcheson
bDepartment of Head and Neck Surgery, Section of Speech Pathology and Audiology, MD Anderson Cancer Center, Houston, TX, USA
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b, Avraham Eisbruch
Avraham Eisbruch
cDepartment of Radiation Oncology, University of Michigan, Ann Arbor, MI, USA
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c, Jonathan J Beitler
Jonathan J Beitler
dDepartments of Radiation Oncology, Otolaryngology and Medical Oncology, Emory University School of Medicine, Atlanta, GA, USA
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d, Johannes A Langendijk
Johannes A Langendijk
eDepartment of Radiation Oncology, University Medical Center Groningen, University of Groningen, Groningen, The Netherlands
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e, Anne W M Lee
Anne W M Lee
fCenter of Clinical Oncology, University of Hong Kong - Shenzhen Hospital, Shenzhen, China
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f, June Corry
June Corry
gDivision of Radiation Oncology, Peter MacCallum Cancer Centre, Melbourne, Victoria, Australia
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g, William M Mendenhall
William M Mendenhall
hDepartment of Radiation Oncology, University of Florida, Gainesville, FL, USA
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h, Robert Smee
Robert Smee
iDepartment of Radiation Oncology, The Prince of Wales Cancer Centre, Sydney, NSW, Australia
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i, Alessandra Rinaldo
Alessandra Rinaldo
jUniversity of Udine School of Medicine, Udine, Italy
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j, Alfio Ferlito
Alfio Ferlito
kCoordinator of the International Head and Neck Scientific Group
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k
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aDepartment of Radiation Oncology, Institute of Oncology, Ljubljana, Slovenia
bDepartment of Head and Neck Surgery, Section of Speech Pathology and Audiology, MD Anderson Cancer Center, Houston, TX, USA
cDepartment of Radiation Oncology, University of Michigan, Ann Arbor, MI, USA
dDepartments of Radiation Oncology, Otolaryngology and Medical Oncology, Emory University School of Medicine, Atlanta, GA, USA
eDepartment of Radiation Oncology, University Medical Center Groningen, University of Groningen, Groningen, The Netherlands
fCenter of Clinical Oncology, University of Hong Kong - Shenzhen Hospital, Shenzhen, China
gDivision of Radiation Oncology, Peter MacCallum Cancer Centre, Melbourne, Victoria, Australia
hDepartment of Radiation Oncology, University of Florida, Gainesville, FL, USA
iDepartment of Radiation Oncology, The Prince of Wales Cancer Centre, Sydney, NSW, Australia
jUniversity of Udine School of Medicine, Udine, Italy
kCoordinator of the International Head and Neck Scientific Group
Address for correspondence: Primož Strojan, MD, PhD, Department of Radiation Oncology, Institute of Oncology Ljubljana, Zaloška 2, SI-1000 Ljubljana, Slovenia, pstrojan@onko-i.si
Issue date 2017 Sep.
PMC Copyright notice
PMCID: PMC5902026 NIHMSID: NIHMS952353 PMID: 28759822
The publisher's version of this article is available at Cancer Treat Rev
Abstract
Radiotherapy (RT) is used to treat approximately 80% of patients with cancer of the head and neck. Despite enormous advances in RT planning and delivery, a significant number of patients will experience radiation-associated toxicities, especially those treated with concurrent systemic agents. Many effective management options are available for acute RT-associated toxicities, but treatment options are much more limited and of variable benefit among patients who develop late sequelae after RT. The adverse impact of developing late tissue damage in irradiated patients may range from bothersome symptoms that negatively affect their quality of life to severe life-threatening complications. In the region of the head and neck, among the most problematic late effects are impaired function of the salivary glands and swallowing apparatus. Other tissues and structures in the region may be at risk, depending mainly on the location of the irradiated tumor relative to the mandible and hearing apparatus. Here, we review the available evidence on the use of different therapeutic strategies to alleviate common late sequelae of RT in head and neck cancer patients, with a focus on the critical assessment of the treatment options for xerostomia, dysphagia, mandibular osteoradionecrosis, trismus, and hearing loss.
Keywords: head and neck cancer, radiotherapy, toxicity, late sequelae, treatment
Introduction
Radiotherapy (RT) plays a pivotal role in the treatment of patients with head and neck cancer (HNC). Approximately 80% (range 73.9 – 84.4%) of all HNC patients will receive RT at least once during the course of their disease .
The main principle of action of RT is to restrict the reproductive potential of tumor cells to induce cell death through apoptosis, necrosis, mitotic catastrophe, senescence, and autophagy . Not only tumor cells are irradiated during an RT course. Some degree of damage is inflicted on normal cells adjacent to the tumor and associated target volumes. The clinical manifestation of resultant damage depends on the radiation sensitivity and cellular organization of the irradiated tissues as well as on the distribution pattern, temporal and geographical, of the radiation dose accumulated in these tissues .
Acute side effects that develop during RT can negatively impact its execution but usually resolve within weeks or months after the completion of RT. Late complications of radiation arise, by conventional definition, 3 months or more after the completion of RT. Many late effects progress over time and can, in the long run, adversely impact the quality of patients’ lives [4–6]. Relative recovery or progression of late effects may vary depending on the nature of the problem. For instance, feeding tube dependence may decrease in initial years of survivorship with adequate follow-up [7–9]. The relationship between the duration and severity of acute side and late effects is also well established. The classic concept of consequential late side effects persisting as adverse outcomes of severe acute side effects not completely repaired has also been supported in the literature, particularly in the area of acute mucositis in relationship to late dysphagia [10,11].
The volume of normal tissues irradiated is determined by the size of the target and RT technique . Modern RT technology allows the effective implementation of several radiation beams, precisely collimated with simultaneous modulation of their intensity, which results in the formation of a high-dose volume that can conform to the three-dimensional shape of the target. Improved conformity and relatively steep dose gradients created on the margin of a high-dose region offers better protection of normal tissues that do not overlap target volumes from higher radiation doses. At the same time, a larger amount of body tissues is irradiated to a lower dose when intensity modulated radiotherapy (IMRT), stereotactic radiosurgery, or other modern RT techniques are used . Similarly, many functionally critical structures, such as pharyngeal constrictors and laryngeal framework, are near target volumes such that complete avoidance of toxicity is not always possible even with advanced, highly conformal RT delivery methods.
No technology can entirely protect normal tissues from irradiation and patients will always experience some degree of radiation-associated toxicity. In the head and neck region, perhaps the most prevalent and challenging adverse late effects of RT are connected with the impaired functioning of the salivary glands and swallowing apparatus . However, other tissues and structures in the region can also be at risk, depending mainly on the location of the tumor, e.g. the mandible and hearing apparatus. In the context of modern treatment scenarios, in many patients the damage caused by radiation is further aggravated by concurrent systemic cytotoxic agents .
The most effective way to limit RT-induced toxicity is to decrease the exposure of functionally critical normal tissues adjacent to the target from excessively high radiation doses that are above the tolerance threshold levels of organs at risk (OAR). Because this is not always possible, either due to limitations of available technology or to the proximity of the target to important normal structure(s), strategies were developed to counteract the existing damage caused by RT. The objective of the present review is to evaluate the available methods used to alleviate the most common late sequelae of RT in HNC patients, i.e. xerostomia, dysphagia, mandibular osteoradionecrosis (ORN), trismus, and hearing loss.
Xerostomia
Xerostomia is the subjective sensation of a dry mouth characterized by a marked decrease and/or thickening of the saliva. It results from salivary gland hypofunction (reduced volume of saliva secretion) or a change in salivary composition and is usually associated with problems with swallowing and speech, and oral health. RTOG grade 3/4 xerostomia is statistically related to excess burden in emotional functioning, fatigue, social functioning and sleeping domains on quality of life (QOL) scales, indicating the multidimensional nature of xerostomia [6,16]. Xerostomia scores may worsen even 5 years or more post-RT .
Stimulated saliva production depends on the major salivary glands. During sleep, the main source of saliva is the submandibular glands, although the contribution of minor salivary glands is not negligible. The content and production of saliva by different salivary glands display circadian variations affecting different aspects of symptoms related to salivary dysfunction [18,19]. Objective presentation of salivary gland hypofunction does not necessarily reflect a subjective perception of xerostomia .
RT is a well-known cause of xerostomia although other factors, i.e. gender, age, smoking, alcohol-containing mouthwashes, concurrent platinum-based chemotherapy, and some medications, also play a role . Among different salivary glands in the head and neck region, minor palatal glands have both greater resistance and recovery when compared with the serous secreting parotid glands .
The reported incidence of moderate to severe late xerostomia in studies implementing two-dimensional RT planning was 60–75% and was around 40% in more recent series employing modern RT techniques, such as IMRT [6,23,24]. In the randomized study of Kam et al. , comparing 2-dimensional RT and IMRT in nasopharyngeal cancer patients, better saliva and observer-reported xerostomia was found in the IMRT group with no difference in patient-reported xerostomia between the two groups. In another randomized controlled trial (PARSPORT), HNC patients were randomly assigned to receive either three-dimensional conformal radiotherapy (3D-CRT) without parotid sparing or IMRT . Grade 4 salivary flow dysfunction decreased from 100% to approximately 50% with IMRT in the PARSPORT trial, while significant reduction were observed for both physician-rated as patient-reported xerostomia in the first 2 years after completion of treatment. The reason for this failure of parotid-sparing IMRT is the need to spare not only the serous parotid glands but also the mucinous submandibular and minor salivary glands in the oral cavity . Thus, despite significant improvements in RT planning and delivery introduced during the last number of decades, the need for effective alleviation of dry mouth syndrome is still present in a substantial proportion of HNC patients.
Pilocarpine
Pilocarpine is a nonselective cholinergic parasympathomimetic agonist which exerts its action by stimulating muscarinic receptors on the surfaces of the salivary gland cells. Thus, parts of the gland with preserved function are compulsory for a positive pilocarpine effect and, from this viewpoint, the pattern of dose distribution is important. It has been shown that stem cells and progenitor cells reside in the region of the parotid gland containing the major ducts and that prophylactic pilocarpine treatment enhances the proliferation of undamaged acinar, and progenitor cells in a rat model [27,28].
Pilocarpine was extensively studied concurrently with RT and its protective effect has recently been confirmed by meta-analysis of 8 randomized phase III studies, though only on unstimulated salivary flow . In post-RT patients who have already developed xerostomia, the use of pilocarpine can also be recommended for the improvement of xerostomia (Table 1). According to the results of two randomized placebo-controlled trials, approximately 50% of treated patients will benefit from reactive use of oral pilocarpine; yet, in placebo groups the response rates were in the range of 25% [31,32]. Symptoms of post-irradiation xerostomia exhibited a time dependent manner of improvement with the best response observed at the end of the treatment period, probably due to changes in the oral mucosa as pilocarpine increased the amount of saliva production. Pilocarpine treatment was found to increase stimulated and unstimulated whole salivary flow rates and those from the parotid glands and the palatal glands [22,31,32]. However, the increase in the salivary flow varied through the treatment period and did not necessarily correlate with the subjective improvement. The latter was observed even after a small increase in the amount of saliva secretion.
Table 1.
Cholinergic agonists for radiotherapy-induced xerostomia in patients with head and neck cancer: randomized-controlled studies
| Study | No. of patients
| Doses & therapy duration | Results (drug vs. placebo) |
:-: :-- |
| Total | Drug | Placebo |
:-:
| Johnson et al., 1993 (31) | 207 | 142 | 65 | Pilocarpine,
5 mg t.i.d. for 12 weeks (73 patients)
or
10 mg t.i.d. for 12 weeks (69 patients) |
improvement in overall condition of xerostomia: 53.5% vs. 42.9% vs. 25%, p=0.010
improvement in saliva production at weeks 4 and 8; did not correlate with symptomatic relief
withdrawal from the study (adverse experience): 5.5% vs. 29% vs. 3.1%
no serious drug-related adverse events
|
| LeVeque et al., 1993 (32) | 162 | 75 | 87 | Pilocarpine,
Weeks 1–4: 2.5 mg t.i.d.
Weeks 5–8: 5 mg t.i.d.
Weeks 9–12: 10 mg t.i.d. |
improvement in overall condition of xerostomia: 48.7% vs. 28.1%, p=0.015
improvement in whole saliva and unstimulated parotid saliva flow
withdrawal from the study (adverse experience, lack of efficacy): 14.67% vs. 13.75%
no serious drug-related adverse events
best results: doses >2.5 mg t.i.d., continuous treatment for 8–12 weeks
|
| Chambers et al., 2007 (39) | | | | |
global xerostomia improvement: 47.4% vs. 33.3%, p=0.0162
improvement in unstimulated saliva flow: p<0.0001 vs. p=0.942
|
| Study 003 | 284 | 139 | 145 | Cevimeline,
30 mg t.i.t. for 12 weeks |
withdrawal from the study (adverse experience): 14.6% vs. 3.5%
treatment-related adverse effects: 41.6% vs. 17.4% (mild to moderate)
|
| Study 004 | 286 | 139 | 147 | If no improvement in dry mouth: escalation of dose to 45 mg t.i.d. at week 6 |
global xerostomia improvement: 48.9% vs. 47.6%, p=0.9565
improvement in unstimulated saliva flow: p=0.0002 vs. p=0.0928
withdrawal from the study (adverse experience): 13.9% vs. 5.5%
treatment-related adverse effects: 38% vs. 18.6% (mild to moderate)
|
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The main problems related to pilocarpine treatment are the fact that the beneficial effect of the drug expires soon after termination of drug administration, and the adverse reactions due to cholinergic stimulation, among others bronchospasm, bradycardia, vasodilatation and diarrhea. In the reported randomized studies, toxicity was common (in a great majority of the patients), usually described as mild but occasionally also of higher grades, and dominated by sweating, flu-like syndrome, and urinary symptoms which generally occurred within 60 min after dosing and were short lived. The adverse effects were the reason for premature withdrawal from the studies in 9–16.9% of the patients [31–33], and were found to be clearly dose-dependent as 5.5% of the patients in the 5-mg group and 29% in the 10-mg group had to stop treatment due to toxicity . Continuous pilocarpine treatment for more than 8 weeks with doses around 5 mg three times a day resulted in the best efficacy/toxicity ratio, with some patients experiencing additional benefit at doses up to 10 mg t.i.d [32,33]. However, due to the irreversible character of the salivary glands’ impairment by irradiation the life-long administration of the drug is required which can be a problem or even contraindicated in some patients with certain pulmonary, cardiovascular or eye diseases, despite the generally low toxicity profile of the drug. According to the results of a maintenance study reported by Jacobs and van der Pas , during 36 months of 5 mg t.i.d pilocarpine treatment, 18.1% of 265 included patients discontinued use of the drug due to an adverse experience but, on the other hand, lack of efficacy was the reason in only 12.8% of cases and the initial drug dose was safely escalated above 5 mg in 48.3% of the patients. There was a significant improvement in all the evaluated criteria of oral function and no evidence of a decrease in therapeutic effect was observed over time.
To alleviate the systemic toxicity of pilocarpine and to maximize the local response other forms of the drug were tested, i.e. a topical administration of pilocarpine suspended in candy-like pastilles , a lozenge , or a mouthwash [37,38]. As these studies were small, further investigations are warranted to confirm some of the promising results observed.
Cevimeline is another cholinergic agonist with a high affinity for muscarinic M3 receptors which has also been tested in a post-irradiation setting. It was found to be effective in a dose of 30–45 mg t.i.d. administered orally for 12 or 52 weeks with a documented increase in unstimulated salivary flow rate, improvement of oral dryness, and mild-to-moderate though frequent side effects [39,40]. In 17.6% of the patients cevimeline-related toxicity was the reason for abandoning the study medication. Recently, cevimeline was compared with oral pilocarpine in a small randomized, cross-over, double blind study, showing a slightly higher, though non-significant, incremental increase in saliva production at the end of four weeks with pilocarpine .
Acupuncture
The stimulation of salivary gland secretion with acupuncture and the alleviation of xerostomia can only be seen when a portion of the salivary gland remains functional. Little is known about acupuncture’s mechanisms that could explain its therapeutic qualities . Regeneration of salivary gland tissue after electrical stimulation of the parasympathetic nerve to the parotid and submandibular glands has been demonstrated in an animal study . True acupuncture was found to be associated with neuronal activations which were absent during sham acupuncture stimulation . Acupuncture was found to be effective also in pilocarpine-resistant patients and as a preventive measure . Recently, an attempt to standardize the procedure in patients with radiation-induced xerostomia was reported .
Acupuncture and acupuncture-like transcutaneous electrical nerve stimulation (ALTENS) were evaluated in 5 clinical randomized trials with modest to high degrees of clinical benefit observed in all trials (Table 2, [48–52]). An increase in whole saliva secretion, and subjective alleviation of related symptoms, was observed for at least 6 months after intervention and up to 3 years with additional acupuncture therapy . The favorable toxicity profile of acupuncture and ALTENS was associated with good treatment tolerance in all reported studies which is of paramount importance when considering a maintenance strategy.
Table 2.
Acupuncture for radiotherapy-induced xerostomia in patients with head and neck cancer: randomized-controlled studies
| Study | No. of patients
| Interventions | Results (intervention vs. control) |
:-: :-- |
| Total | Groups |
:-: |
| Blom et al., 1996 (48) | 38 | 20 | Real AP (2 series of 12 THs with 2-week pause 2 THs/week, 20 min/AP |
improvement in saliva flow rates: in 68% vs. 50% patients
no significant difference between real AP and sham AP
|
| 18 | Control: superficial (sham) AP |
improvement persists during the observation year in both groups
control group: improvement is smaller and appears later
|
| Cho et al., 2008 (49) | 12 | 6 | Real AP (2 THs/week, 20 min/TH, for 6 weeks) |
real AP significantly increase unstimulated salivary flow rate and improve the score for dry mouth (at 6 weeks)
|
| 6 | Control: superficial AP |
| Pfister et al., 2010 (50) | 58 | 28 | AP (1 TH/week, 30 min/TH, for 4 weeks) |
real AP produces a significantly greater improvement in reported
|
| 30 | Control: usual care |
xerostomia compared to “usual care”
|
| Simcock et al., 2013 (51) | 145 | 70 | AP → oral care |
no significant changes in stimulated or unstimulated saliva flow by AP
|
| 75 | Oral care → AP |
| AP: 1 TH/week, 20 min/TH, for 8 weeks |
AP provides significantly better reductions in patient reports of xerostomia-related symptoms compared to “oral care”
|
| Oral care: 1 educational sessions of 1h, 1/month |
| Wong et al., 2015 (52)1 | 148 | 75 | ALTENS (2 TH/week, 20 min/TH, for 12 weeks) |
no significant difference in the change in xerostomia burden between the ALTENS a pilocarpine
|
| 73 | Pilocarpine (5 mg t.i.d., for 12 weeks) |
significantly less toxicity in ALTENS group
|
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AP – Acupuncture; TH – Therapy; ALTENS – Acupuncture-like transcutaneous electrical nerve stimulation.
1Only 96 and 76 patients with all items completed were evaluable at 9 and 15 months from randomization.
In two prospective randomized, controlled trials classical (real) acupuncture was compared with superficial (placebo/sham) acupuncture [48,49] and in two other randomized trials with usual oral care or group oral care education (ARIX trial, ). Although the sample sizes were small (12–58 patients), with one exception (145 patients; ), subjective benefits were recorded from acupuncture in all of these trials using different measurement tools. However, no objective gain from the intervention was observed as the flow rates, though increased over time by acupuncture, also improved in the control groups, which resulted in non-significant differences between the experimental and control groups [48–50]. Positive outcomes observed in the control groups may be due to the stimulation of skin receptors by the superficial needle insertion as a placebo, which also enhance the functional recovery of the salivary glands, though to a lesser degree (sham acupuncture); a therapist-patient relationship (i.e. the Hawthorne effect); or the patient’s expectation of benefit from the intervention. Only the ARIX trial controlled for these effects by a randomized crossover design, giving both interventions (acupuncture and oral care educational sessions) in both arms; yet, acupuncture resulted in improved xerostomia-related QOL compared to educational sessions .
In a recently reported RTOG 0537 phase III randomized study, ALTENS was compared with standard pilocarpine treatment (5 mg t.i.d. for 12 weeks); there was no difference in radiation-induced xerostomia symptom burden at 6 months after treatment between the two arms . However, there was a consistent trend towards greater improvement in the xerostomia quality of life scale (XeQOLS) scores at all follow-up time points, a higher percentage of positive responders (defined as ≥20% improvement from the baseline XeQOLS scores; at 12 months’ post-therapy, 83% vs. 63%, p=0.04), and significantly less non-hematological toxicity (CTCAE v3.0 grade ≤3, 20.8% vs. 61.6%) in patients receiving ALTENS. Notably, the statistical power of the trial was limited due to the low proportion of evaluable patients (52 of 148, 35.1% of recruited patients failed to complete the planned assessments, particularly in the pilocarpine group).
Other strategies and emerging approaches
These include hyperbaric oxygen therapy (HBOT), gustatory/masticatory stimulation and the use of lubricants or saliva substitutes, whereas gene therapy and stem cell therapy are recent emerging approaches.
The assumption that HBOT may provide some relief to patients suffering from RT-induced xerostomia was based on its effectiveness in the treatment of bone ORN and soft tissue necrosis. Only a few studies analyzed the role of HBOT in post-RT xerostomia, with a limited number of patients included, and were extensively elaborated in a systematic review by Fox et al. . In six studies, 227 patients were treated 2 years or more after RT; all but one study were prospective cohort trials and one was a randomized controlled trial, with an average number of HBOT dives between 20 and 43, respectively. [56–62]. Increases in the stimulated salivary flow rates [60,61] and scores on the visual analog scale quantifying xerostomia severity [58,61], as well as an improvement in the salivary viscosity and the EORTC H&N35 scores for sensation of dry mouth [56,57,62] were demonstrated but not accompanied by an overall QOL improvement within a year of completion of HBOT [57,62]. Some studies suggested a long-term benefit of HBOT, up to 18 months after its completion. However, the low number of treated patients, only one study with a control group (i.e. not allowing a fair discrimination between the HBOT effect and the placebo effect), lack of appropriate details on RT reported, heterogeneous patient populations (regarding indications for HBOT as the majority was treated for other RT sequelae which could confound the results), and the possibility of publication bias (negative studies may have not been published) limit the strength of the conclusions. Also, no risk assessment was done for cancer recurrence .
The utilization of gustatory or masticatory stimulants, lubricants and saliva substitutes is a purely symptomatic measure. Nevertheless, it could be an esteemed provision in patients with insufficient residual secretion of saliva on stimulation. These agents cannot replace the antimicrobial and immunological properties of saliva. A wide variety of different acid-tasting substances, including sugar-free chewing gum, sweets, vegetables and fruits, and pharmacological sialologues have been tested in small cohort studies but also randomized controlled trials with different levels of success. In general, gustatory and masticatory stimulation by acidic substances showed some increase in whole saliva secretion and amelioration of oral dryness; oral lubricants and saliva substitutes usually exert a short-term effect over the placebo effect . Constituents resembling the physical properties of glycoproteins and the antibacterial components of saliva are commercially available in the form of over-the-counter gels, mouthwashes, or sprays. Chemically, they are based on carboxymthylcellulose, mucin or xanthan gum, the latter two being characterized with superior rheological and wetting properties and, therefore, preferred by patients with xerostomia [64–67]. Recommendations for the use of oral lubricants and saliva substitutes to alleviate hyposalivation as proposed by Regelink et al. are shown in table 3.
Table 3.
Recommendations for the use of oral lubricants and saliva substitutes to alleviate hyposalivation (adapted from Regelink et al., Ref.68)
| Xerostomia | Recommendation |
:-- |
| Severe | During the night: |
| saliva substitute with gel-like properties or mucin-containing lozenges |
| During the day: |
| saliva substitute with properties resembling the viscoelasticity of natural saliva (e.g. polyacrylic acid, xanthan gum, or mucin based substitutes) |
|
|
| Moderate | During the night: |
| gel |
| During the day |
| saliva substitute with a rather low viscoelasticity (e.g. carboxymethylcellulose, hydoxypropylmethylcellulose, mucin-based sububstitutes); or low concentrations of xanthan gum and polyacrylic acid |
|
|
| Slight | Little amelioration is expected from the use of saliva substitutes (gustatory or pharmacologic stimulation is the treatment of choice) |
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Gene therapy is based on a viral vector injection, conveying genetic information into a tissue to result in some beneficial change. It represents a promising new approach to the treatment of RT-induced xerostomia. Salivary gland tissue was found to be a good target for gene transfer due to stable, well-differentiated and slowly dividing epithelial cells, limited leakage of the vector, and direct access via duct orifices . The most studied is the adenovirus transfer of the human aquaporin-1 gene (hAQP1), which encodes a water channel protein involved in the osmotic movement of water in radiation-surviving salivary duct epithelial cells in the damaged gland. The strategy proved to be safe and effective in a phase I/II clinical study in humans previously irradiated for HNC, showing an increase in the saliva flow rate and a reduction of xerostomia-related symptoms which can continue years after hAQP1 delivery [70,71]. Recently, a non-viral approach employing ultrasound-assisted AQP1 cDNAs transfer was successfully tested in an animal model, allowing multiple gene administrations to maintain elevated salivary flow as opposed to the viral vectors where only a single administration is possible . The hedgehog pathway appeared to be another gene therapy target to overcome RT-induced hyposalivation. Its transient activation in mouse salivary glands by overexpressing the Sonic hedgehog transgene or administering a smoothened agonist showed the potential to restore salivary function in pre-irradiated glands [73,74].
Like other organs, also in salivary glands, the stem cells play a central role in tissue homeostasis . As precursors of progenitor cells, which have a lower self-renewal capacity and may be tissue-specific, they focus on extensive tissue regeneration and salivary bioengineering research at the cell culture level, and in animal models but not yet in humans. Stem/progenitor cells from the major and minor salivary glands have been successfully isolated, cultured and transferred to rat glands , though minor salivary glands, bone marrow and adipose tissue-derived mesenchymal stem cells have also been identified as a possible source for the repair of radiation-damaged tissue [77–79]. In addition, the self-duplication potential of already differentiated acinar cells, which also play a role in salivary gland homeostasis (i.e. the replacement of aging and injured cells), appears to be another option to counteract salivary gland dysfunction and xerostomia .
Current clinical standards
Despite substantial progress with various approaches tested in clinic to alleviate RT-induced xerostomia, there are no consensus standards for xerostomia management. No uniform treatment protocol can be recommended for wide clinical adoption, and still limited advice can be given to the patient who is complaining over a dry mouth. The empirical first step is to introduce symptomatic measures recommending gustatory and masticatory stimulants, lubricants and saliva substitutes. However, many patients may prefer frequent moistening of oral tissues with water rather than the use of saliva substitutes due to the short duration of relief they provide. A published recommendation is to test different saliva substitutes by individual patients to select the most effective one as her/his preference may also be of importance for the success of this treatment . In patients without history of cardiovascular and eye diseases (glaucoma, inflammation or infection) or uncontrolled asthma, a course of pilocarpine 5 mg t.i.d can be tested with possible increase to 10 mg t.i.d in selected patients who do not experience adverse effects. The duration of pilocarpine therapy should be adjusted to the drug effectiveness and exerted side effects. If available, acupuncture is an option in those not amenable for or resistant to pilocarpine therapy.
Dysphagia
Difficulty in swallowing is a common consequence of RT for HNC. In a prospective cohort study of 238 HNC patients treated with 3-dimensional conformal RT (65%) or IMRT (35%) concurrent with chemotherapy (11%), the prevalence of grade 2–4 swallowing dysfunction according to the RTOG/EORTC Late Radiation Morbidity Scoring Criteria was 22% at 6 months, and 14% at 12 and 24 months . Late dysphagia was recognized as one of the key adverse factors in QOL testing of HNC patients, with feeding-tube dependency having the most negative impact [80,81]. Similarly, clinician-graded dysphagia correlates strongly with QOL with larger effect sizes than xerostomia in long-term follow-up after curative RT . Fortunately, permanent feeding tube dependency in disease-free patients, particularly those treated with RT alone, is rare [7,83]. However, even mild late dysphagia is a major correlate of QOL . Moreover, it is well recognized that absence of a feeding tube does not equate to a safe or efficient swallow. For instance, half of radiographically-detected chronic aspirators after larynx preserving RT are gastrostomy free, essentially electing to eat despite their dysphagia but with excess risk of pneumonia and related morbidity .
Dysphagia is a complex neuromuscular toxicity. More than 25 pairs of muscles, located from the oral cavity to the esophagus, are involved in swallowing. Swallowing requires extreme precision to coordinate the hemispheric cortical centers, the brainstem central pattern generator, associated cranial nerves, muscles, and sensory receptors . Impairment in any segment of this complex process can result in an aberrant cough reflex and aspiration with consequent pneumonia and chronic bronchial inflammation as aspects of dysphagia alone or together with oropharyngeal inefficiency. Inefficient bolus clearance impacts risk of prolonged or even permanent feeding tube dependency, weight loss, nutritional deficiencies and life-threatening malnutrition .
Radiation-associated dysphagia (RAD) evolves in various ways as recently demonstrated with longitudinal cluster analysis of 238 consecutive patients with HNC followed serially with clinical assessment of dysphagia for 24 months . Using clinician graded RTOG/EORTC dysphagia endpoint, roughly 16% of patients treated with definitive 3D-CRT or IMRT had substantial late burden of dysphagia attributable to one of two patterns: “severe persistent” [8%, grade ≥2 at 6 months that remained up to 2 years) or “progressive” [8%, grade <2 at 6 months, subsequently progressed to grade ≥2) dysphagia. Distinct patterns of RAD were felt to possibly signal different biologic processes underlying swallowing dysfunction in long-term survivors.
Mechanisms of late dysphagia after RT for HNC are poorly understood. Soft tissue fibrosis has long been considered the primary source of RAD, with associated restriction in the compliance and contractility of underlying musculature due to post-inflammatory scarring processes and lymphedema. Compounding the dysfunction caused by local RT damage, muscle atrophy (with associated weakness) may also result from disuse of the oropharyngeal musculature during RT when patients often stop eating normal foods and may require several months of tube feeding to sustain nourishment while acute RT toxicities are at their peak . Sensory loss is quite poorly understood, yet is likely another underreported contributor to RAD, given that roughly half of chronic aspirators do so silently without the normal sensory response to clear the airway of foreign bolus entry [88–90].
Baseline function and early recovery of swallowing are critical clinical predictors of late effects. Pre-treatment swallowing dysfunction induced by the tumor-driven destruction of normal tissue is possible, particularly with locally advanced stage disease. Baseline dysfunction increases risk of chronic or persistent RAD . Likewise, the prognostic relevance of acute radiation-induced symptoms for the development of late dysphagia underscores the etiologic role of inflammation . Edema is currently under study in a prospective cohort as another acute sequela along the continuum to soft tissue fibrosis, suggesting the potential adjunctive role of lymphedema therapy strategies for dysphagia management . The etiologic role of inflammation and edema may be more relevant for classic forms of persistent dysphagia as a consequential late effect of acute reactions. There is emerging work to suggest a unique subtype of dysphagia in HNC survivors referred to as late- RAD [93,94] Late-RAD is defined by presentation with severe oropharyngeal dysphagia after a long interval of adequate functioning after acute effects of RT resolved. Early work suggests a median latency of 8 years to presentation of late-RAD with de novo lower cranial neuropathies preceding the deterioration in swallowing function in most cases. Thus, denervation appears to be a major source of the neuromuscular dysfunction in late-RAD rather than predominant edema and stricture as in more classic, earlier forms of RAD.
The degree of swallowing dysfunction depends on the distribution of the delivered RT dose. Thus, substantial efforts are underway to understand how best to optimize radiation planning to achieve better and more durable swallow preservation. Pharyngeal constrictors and the larynx were first identified by multiple investigators as primary dysphagia-aspiration-related organs (DARS) . Beyond the classic DARS, emerging work suggests relevance of dose/volume distributions to submental musculature (i.e., geniohyoid, mylohyoid) that are critical to elevate the hyolaryngeal complex and assist with cervical esophageal opening [96,97]. Reporting normal tissue complications probabilities as a function of DARS, Eisbruch et al. gave solid ground for efforts to reduce dysphagia in the IMRT era [95,98,99]. Work by Christianen et al. further supports the dose constraints to classic DARS, finding mean dose to the superior pharyngeal constrictor muscle and supraglottic larynx as the most important predictive parameters for the occurrence of grade 2–4 RTOG/EORTC swallowing dysfunction at 6 months after completion of (chemo)RT. Interestingly, the importance of the latter (supraglottic dose) declined over time, while the effect of the dose to the superior constrictor remained stable, suggesting time dependence of these parameters [9,100]. In certain clinical situations, e.g. in patients with deep posterior wall pharyngeal primaries or in those with extensive base of tongue and supraglottic tumors, injury to the important swallowing structures with resulting high-grade dysphagia is unavoidable, diminishing the potentially positive effects of prophylactic measures for improving swallowing outcome in these subgroups of patients . The risk of tube feeding dependence was found to depend on a variety of dose-volume histogram (DVH) parameters, including the superior and inferior pharyngeal constrictors and the contralateral parotid gland mean dose . Recently, Christianen, et al. showed that when the dose to the superior pharyngeal constrictor and supraglottic area was reduced by adding dose constraints to these structures, the risk of grade II–IV dysphagia could be reduced significantly as well.
Behavioral and exercise-based swallowing rehabilitation
The exercises, maneuvers and postures aimed at improving swallowing efficiency and safety include motor exercises (to increase the strength, mobility, and endurance of the swallowing structures); patterning and postural patterning and techniques (to minimize aspiration/maximize swallowing efficiency with changes in body position or motor patterns of swallowing); and sensory techniques (to increase sensory response by stimulation, employing alterations in temperature, taste, pressure, or bolus – its consistency, placement, or size) . Neuromuscular electrical stimulation coupling surface electrodes placed overlying neck or submental muscles during exercise and respiratory-swallowing training using visual feedback for the optimization of respiratory/swallowing patterning also belong to this group of behavioral interventions, often directed by speech pathologists for patients in North America [105,106].
Precious few studies have examined rehabilitation interventions introduced after HNC (chemo)RT (with or without surgery) to alleviate dysphagia [104–117]. In six randomized studies, two different swallowing strategies were compared rather than to the non-interventional control group [111–117]. While there is single-institution evidence supporting modest effect achieved after non-preventive behavioral or exercise based swallowing interventions after (chemo)RT, they are of limited quality . Major and recurrent limitations include the small sample size and methodological issues, including a lack of standardization in the evaluation methods and the absence of a core set of dysphagia parameters for reporting, and almost universal short-term follow-up of patients. These issues are summarized in a recent inconclusive Cochrane review on the topic .
On the contrary, a recently published double–blinded, randomized controlled trial comparing a combination of neuromuscular electrical stimulation and swallowing exercises to exercise with sham stimulation in 170 post-RT dysphagia patients reported no benefit of electrical stimulation, and swallowing exercises alone (with sham stimulation) were found to have modest effect on diet and QOL and no improvement on radiographic measures of swallowing function [116,117]. The authors concluded that traditional swallowing exercises may not be effective in HNC patients with moderate-to severe chronic RAD, perhaps due to progressive fibrotic changes and entrapment of the swallowing apparatus induced by RT. Contrary to these disappointing objective results, the patients from the trial reported significantly better diet and QOL outcomes, as is often the case with conflicting patient-reported and physical measures. In RTOG 9003 study, the 1- and 5-year feeding tube rates in disease-free patients were 12.1% and 7.8%, reflecting that patients can adapt and force themselves to some extent to eat with impaired pharynx .
Novel, systematic efforts are under study to offer more for complex dysphagia. These include bolus-driven “boot camp” style therapies such as McNeil Dysphagia Therapy Program and biofeedback swallow pattern training [120–123]. These therapy models focus on standardized, clinician-directed training to adjust the swallow pattern during bolus swallows rather than focusing on strengthening exercise outside of mealtime. Small, single institution case series report functional gains (better diet or less frequent aspiration) but persistent physiologic dysphagia on videofluoroscopy that can be interpreted at helping compensation rather than reversing dysphagia [120–123]. These results, together with perceived non-uniformity in usual dysphagia therapy practices, call for the development and refinement of novel dysphagia interventions tailored to unique pathophysiology in HNC patients on one hand and to the development of evidence-based algorithms for personalized HNC dysphagia treatment guidelines [124–126].
Acupuncture
While acupuncture appears have some efficacy as a therapy for post-stroke dysphagia, its role in RAD is yet to be determined . Although the etiology of dysphagia in stroke and HNC differ, the two share many aspects of the needling technique . In the literature, 3 reports were published, 2 of them from China, that evaluated the efficacy of acupuncture in HNC patients with dysphagia [129–131]. In the Dana-Faber Cancer Institute study, a subjective improvement of various degrees in patient-reported swallowing functions, xerostomia, pain, and fatigue levels was reported in 9 of 10 treated patients . Of 7 PEG tube-dependent patients, 6 had their feeding tube removed after acupuncture. Building on these observations, a pilot randomized sham-controlled trial has been conducted, showing that both active acupuncture and sham acupuncture during and after head and neck RT are feasible and safe [128,132]. An improvement in the dysphagia-related QOL from the baseline to 12 months post-RT was observed in both treatment arms, though with no difference between the two.
Current clinical standards
Level 1 evidence is lacking as a basis for evidence-based management of late dysphagia. Late dysphagia is complex and a comprehensive assessment is the first step in personalized management. Assessment ideally should include imaging studies (videofluoroscopy and/or endoscopy) with consideration of rapidly emerging technologies like pharyngeal high resolution manometry to characterize pathophysiology and pattern of dysphagia. Modest effect sizes of current strategies may reflect low or static intensity of therapy programs where a patient practices the same tasks at the same intensity on their own rather than varying performance with progressive resistance or progressive complexity to enhance gains. Personalized, intensive therapies are likely ideal. Published work also suggests impetus to avoid pharyngeal disuse that likely hastens or exacerbates the problem . Preventive efforts are paramount that should couple dose-optimization strategies with supportive care to help patients avoid pharyngeal disuse. Despite the uncertainty of the evidence, patients with dysphagia might still be encouraged to practice swallowing exercises as the inhibitory potential of such exercises on the progression of swallowing dysfunction cannot be ruled out. An additional justification for maintenance exercise regimens draws from the possible positive effect of the simple act of practicing swallowing over and above the meal time effort which could impact the patient’s skill, ease and rate of eating [116,117]. Emerging methods such as acupuncture and manual therapies are much less integrated in routine practice, but hold promise in small clinical studies. Novel strategies and interdisciplinary trials are desperately needed to improve options for late dysphagia in HNC.
Mandible osteoradionecrosis
ORN is defined as a necrotic process in the bone that results from a high-dose RT and persists for 3 months or longer, worsens slowly and does not heal spontaneously . In addition to dosimetric factors, poor dental health before RT, post-radiation extraction of teeth and trauma (by prosthesis, bone biopsy, salvage surgery) in bone regions exposed to a high RT dose, are the main risk factors for ORN development. Xerostomia and trismus are well recognized contributing factors, as they impair oral hygiene and promote caries [133,134]. ORN presents clinically as a painful and denuded bony region, with purulent drainage and sometimes progresses to fistula formation (to the mucosal or skin surface). On a panoramic radiograph, loss of trabecular architecture usually seen though the appearance of the affected bone may be normal. In a historical RT series, the incidence of ORN ranged from 2% to 22%: this wide range reflects not only the differences in RT techniques employed but also ambiguity in definition and difficulties in the diagnosis of ORN [133,134]. In an extensive review of 31 studies published between 1990 and 2008 by Peterson et al. , the weighted prevalence for ORN in patients with different types of HNC treated with conventional RT was 7.4%, with IMRT 5.2%, with (chemo)RT 6.8%, and with brachytherapy 5.3%. Recently, De Felice et al. reviewed 10 studies with over 3,000 patients included, employing exclusively IMRT: the ORN incidence rates ranged from 0% to 6.3% (median 1.4%). In a detailed analysis of 531 oral cavity, oropharyngeal and salivary gland cancer patients treated with IMRT and doses to the mandible ≥60 Gy, Studer et al. reported the overall incidence of ORN as 7%; in oral cavity cancer patients with mandibular surgery the rate was 29% (when no mandibular surgery was done 7%) and when marginal or periostal bone was resected it was 39% (segmental resection or no resection, 7%). The importance of reduced mandibular volumes receiving high doses and improved salivary flow with associated improved oral health when IMRT was used, together with meticulous dental prophylactic care, was stressed by Ben-David et al. , who reported no case of ORN among 176 patients after IMRT.
Available therapeutic options in ORN are complementary to its complex pathophysiology and both currently accepted theories of ORN development. The Marx’s “3Hs paradigm” or “hypovascular-hypoxic-hypocellular” theory, with persistent hypoxia as one of the main features of ORN, provides the grounds for the use of HBOT . Later, Delanian and Lefaix introduced the “fibroantropic” theory, proposing the activation and dysregulation of fibroblast activity and resulting fibrosis of tissue that is prone to traumatic breakdown as the main mechanism of ORN. Several drugs that counteract various aspects of the latter theory have been studied .
The range of possible intervention in patient with ORN starts from non-invasive measures with optimization of oral hygiene and antibiotic coverage, and escalates through sequestrectomy and debridement to major and more mutilating surgical procedures. Surgery employing resection of involved bone and free-flap reconstructions are indicated in advanced-stage ORN cases with fractures and/or fistulas [133,134]. According to a systematic review of the pertinent literature by Lee et al. , the fibula is the workhorse free flap for reconstruction in mandibular ORN, with the risk of flap failure and postoperative complications being significantly increased compared to reconstruction of the primary tumor in unirradiated patients. Thus, it seems appropriate to defer resection and reconstruction in patients whose symptoms can be managed without an operation. A conservative approach also gained confirmation in health-related QOL analyses where the scores were disappointingly poor after mandibular surgery [142,143].
Hyperbaric oxygen therapy
Inhalation of pure oxygen (100%) at pressures greater than 1 atmosphere absolute in an airtight chamber promotes healing by stimulating angiogenesis, epithelization, osteoblastic proliferation, collagen synthesis, and has antibacterial properties. In healthy tissues, it results in hyperoxic vasoconstriction which is followed by the redistribution of peripheral blood volume toward the hypoxic tissues (the Robin-Hood effect) . Lung, brain, and ocular (progression of cataract, transitory myopia) toxicities may occur due to hyperoxia and increased air pressure but are rarely observed when HBOT follows accepted protocols .
In 2010, Peterson et al. reviewed seven studies employing HBOT, alone or with a varying degree of surgical management, for the treatment of ORN. ORN resolution rates in these studies varied from 19% to 93% showing no advantage of HBOT over surgery but rather a synergistic effect when both therapies were combined. This conclusion was recently confirmed in a series of 27 patients reported by Dieleman et al. who observed a beneficial effect of HBOT alone in early-stage disease (improved/healed cases: stage I – 54%, stage II – 25%) but in stage III ORN cases healing was achieved only after a HBOT-surgery combination. In the latter group, HBOT seemed to be ineffective for bone healing but may have resulted in improvement of the surrounding soft tissues which positively affected the healing of the surgical wound. In addition, a better prospect for successful healing after tooth extraction (compared to antibiotic cover) and for complete mucosal cover of exposed bone in patients with ORN when using HBOT was confirmed in Cochrane’s review, though the quality of evidence was only found to be moderate . The only randomized trial was reported by Annane et al. who compared HBOT with a placebo, with (more advanced cases) or without (less advanced cases) surgery. The trial was prematurely stopped after enrolling 68 patients with overt ORN due to a potentially worse outcome in the HBOT arm. This study was severely criticized due to several methodological and performance issues . Wide variability in HBOT regimens used in the UK and Europe concerning the actual pressure, duration and number of dives was reported [148,149]. To overcome these shortcomings, the role of HBOT in the prevention and treatment of ORN is currently being investigated in three prospective multicenter randomized trials: the UK’s HOPON trial, the Danish DAHNCA21 trial, and the Portuguese phase II trial on HBOT therapy with or without antifibrotic agents [150,151].
Medical management
The following agents that cause a blockade at specific points/levels in the fibroantropic model of ORN development were used in the management of ORN: pentoxifylline, tocopherol and clodronate. Pentoxifylline is a methylxanthine derivative with an in vivo established inhibitory effect on fibroblast proliferation and extracellular matrix production, as well as collagenase activity stimulation. Due to toxicity induced by the drug concentration needed for the effective suppression of collagen synthesis, the isolated administration of pentoxifylline is not indicated. Antioxidant tocopherol (vitamin E) is the free radial scavenger, acting as a protector of cell membranes against lipid peroxidation; its antifibrotic activity is based on downregulation of procollagen gene expression. Both agents exhibit anti-inflammatory actions by inhibiting TNF-α,TGF-β and other mediators. The third agent, clodronate, is a first-generation bisphosphonate that inhibits bone resorption by reducing osteoclast activity through direct activation of osteoblasts it increases bone synthesis and decreases proliferation of fibroblast .
Promising early results on the use of a pentoxifylline-tocopherol combination in a mixed patient population with symptomatic RT-induced fibrosis resulted in the phase II trial in patients with refractory mandible ORN . A total of 18 patients were treated with a daily oral combination of 800 mg of pentoxifylline and 1000 IU of tocopherol for 6 to 24 months; the worst 8 cases also received clodronate 1600 mg/day, 5 days a week. Treatment was well tolerated and complete recovery was observed in a median time of 6 months. The favorable clinical outcome and the toxicity profile of the medical ORN treatment were duplicated by the same group in a larger prospective cohort of 54 HNC patients with refractory mandibular ORN, mainly after surgery and HBOT , and was confirmed by several other groups [152,156,157]. However, the true value of this strategy is yet to be confirmed in a randomized clinical trial setting .
Current clinical standards
Currently, there is no gold standard treatment of ORN and no widely accepted guidelines exist. It seems that maximal benefit is gained by combining several therapeutic strategies directed against different targets in the ORN pathophysiology chain, taking into consideration the stage of the condition and characteristic of the patient and her/his malignant disease. Whereas early-stage ORN in patients who are symptom-free or mildly symptomatic needs only conservative measures (a wait-and-see policy: optimal oral hygiene, antibiotics), any sign of progression requires early surgical intervention in the form of sequestrectomy and debridement with local mucosal flaps to cover bone defect. The role of HBOT and medical management, however, is yet to be defined with consensus guidelines and robust clinical trials.
Trismus
RT-associated trismus infers fibrotic changes with contracture in mastication structures, including the masseter and pterygoid muscles, damage to their neural innervation and temporomandibular joint degeneration . These changes can result in a considerable reduction in mouth opening (with a maximum inter-incisor opening <35 mm, i.e. the distance usually used as the cutoff point) compared to pre-irradiation status with about two-thirds of the total reduction in mouth opening observed in the first 9 months after RT . Then, the process becomes slower but may continue over later years. Dosimetric studies showed an increase in the probability of trismus of 24% per every additional 10 Gy delivered to the pterygoid muscles, with doses as low as 15 Gy resulting in functional impairment [160,161]. The prevalence of trismus primarily depends on tumor site and size, being the highest in patients with tumors close to the mastication apparatus, i.e. parotid lesions, nasopharyngeal and lateralized oropharyngeal or posterior oral cavity primaries [162–164]. The wide prevalence range reported in the literature reflects the differences in the studied populations and, consequently, the dose burden across the studies. With improved control over spatial radiation dose distribution, IMRT may lessen the problem with trismus, though not in cases with a target lesion sited in or next to structures important for mastication. In a systematic review of 22 studies, Bensadoun et al. reported the weighted prevalence of trismus following RT of 25.4% for conventional RT, 5% for IMRT and 30.7% for a combination of RT (mainly conventional) and chemotherapy.
Trismus was recognized as the third most burdening side-effect of oncologic therapy . It was found to be associated with poorer health-related QOL and a greater proportion of HNC patients with trismus displayed depression compared to a control group without trismus [164,167]. Improper mastication mandates changes in food consistency and, in extreme cases, enteral tube feeding. Other negative consequences of trismus include difficulties maintaining oral hygiene with an associated increase in the risk of oral infections and dental problems, as well as impaired oral examination, dental care and intubation.
Conservative measures with HBOT and pentoxifylline did not show any convincing improvement of trismus [168,169]. The same was observed for injections of botulinum toxin into the mastication muscles, though they can effectively reduce local pain produced by muscle spasms associated with trismus . Forced mouth opening under general anesthesia can improve trismus, but the effect is usually short-lived and the procedure is less controlled and at considerable risk of fracture and adjacent soft tissue rupture.
Exercise therapy employing different stretching techniques and devices is aimed at increasing the range of mouth opening by strengthening musculature, improving the mobility, flexibility, and elasticity of the temporomandibular joint and by improving blood circulation . Recently, a systematic review by Kamstra et al. identified 12 studies, ranging from chart reviews to randomized controlled trials that analyzed the effect of therapeutic measures on mouth opening after trismus occurred. The duration of the therapy in these studies ranged from 1 to 9 months, the number of exercise sessions per day varied from 2 to 10, with 3 to 8 repetitions, and the duration of stretch ranged from 6 seconds to 60 minutes. Meta-analysis was not possible due to significant heterogeneity in the clinical profile of the patients included, variety mechanical devices used (dynamic bite openers, a sledge hammer device, the TheraBite, Engstron Jaw mobilizing device, tongue depressors, rubber plugs, Dynasplint Trismus System), and inconsistent methodology. Collectively, the studies suggested some efficacy of therapeutic jaw mobility interventions. An increase in mouth opening was reported in the majority of studies, although, in many, the final mean mouth opening was still <35 mm suggesting that even after therapy, a proportion of patients still suffer from trismus. No exercise technique was found to be superior to others, and the results were influenced by compliance with the exercises and the time interval between oncologic treatment and the start of the exercises (the sooner, the better) . A recently reported study using the Dynasplint Trismus System confirmed that early detection of trismus and the early start of exercise therapy ensure a better outcome of mouth opening .
Surgical treatment is an option for patients without known malignant disease with trismus refractory to physical therapy. Coronoidectomy proved effective, increasing inter-incisor distance at least 20 mm, in a series of 18 previously irradiated patients reported by Bhrany et al. . Moreover, all maintained an inter-incisor distance equal to or greater than 35 mm for 6–12 months after the procedure. Similarly, Mardini et al. used free flaps to reconstruct the defects created after surgical trismus release in 11 patients with previous intraoral reconstruction (8 patients had postoperative RT). The initial mean inter-incisor distance of 3.1 mm was increased to 33.4 mm immediately after the release and to 18.9 mm at a mean follow-up time of 22.7 months. The authors concluded that the procedure is a viable option in complex cases that yields a reasonable, long-lasting improvement in mouth opening, intraoral hygiene, and QOL.
Current clinical standard is exercise based therapy with the precise device and schedule still debated. Surgical management is reserved for extreme cases, and HBOT and pharmacological therapy have no strong evidence in published literature.
Hearing loss
Damage to the hearing apparatus with resultant sensorineural hearing loss (SNHL) is a common adverse event after RT [175–177], occurring in up to 43% of irradiated HNC patients. Its incidence is increased to 17–88% when RT is combined with cisplatin . Traditionally, SNHL denotes a pure-tone audiometry confirmed, clinically significant increase in bone conduction threshold at the key human speech frequencies (0.5–4.0 kHz) . Radiation-induced SNHL results from damage to the organ of Corti and stria vascularis of the cochlea, spiral ganglion and/or the acoustic nerve. The pathophysiology includes a small vessel endothelial reaction with vascular insufficiency and insult to the inner ear sensory structures, leading to their progressive degeneration and atrophy, fibrosis and even ossification of the inner ear fluid space. Furthermore, the cochlear nerve can be affected by edema and inflammation in the narrow space of the internal auditory bony canal . Hearing deterioration is worse at higher frequencies compared to speech frequencies and commences soon after RT. Early changes in hearing could be reversible, though the probability of clinically relevant threshold deterioration of hearing increases over time, with the median latency period being 1.5 to 2 years [181,182]. In addition to a radiation dose of 47 Gy or more to the cochlea, the main factors affecting the risk of SNHL are the patient’s age and baseline hearing level, cisplatin dose, post-RT otitis media, and follow-up time [178,179,183–185]. In addition, the deterioration in air conduction caused by RT-induced Eustachian tube dysfunction and middle air fibrosis may accompany persistent SNHL, indicating mixed hearing impairment .
There is no standard treatment of SNHL caused by RT. Different interventions were tested in uncontrolled and retrospective case series but not in a properly designed randomized trial. Conservative measures may include systemic steroids to improve inflammation and edema in the inner ear after radiation-induced damage and to reduce compression of the acoustic nerve due to inflammatory swelling. Chen et al. reported that the use of systemic methylprednisolone during RT can reduce early SNHL caused by irradiation. In a post-RT setting, Sakamoto et al. observed an improvement in hearing loss with oral prednisone, at least in young patients having a useful pretreatment hearing level, if the steroid treatment was administered immediately after the first detection of the hearing loss. The benefit of oral corticosteroid therapy in salvaging acute hearing deterioration was also suggested by Kim et al. . In idiopathic sudden SNHL, however, the value of steroids was termed unclear in the Cochrane review of the evidence , though salvage intratympanic therapy has shown some positive effects in refractory cases .
The rationale for the use of HBOT is based on the role of vascular damage and ischemia in RT-induced SNHL development. A Cochrane review suggests that hearing improvement of 25% can be expected in one out of every five treated patients with acute idiopathic sudden SNHL but no beneficial effect of HBOT was found in chronic cases . In another review by the same source, no evidence of any important clinical effect of HBOT was found on neurological tissues, either peripheral or central .
In patients with RT-induced bilateral profound SNHL, cochlear implants are a viable rehabilitation option. As determined in nasopharyngeal cancer survivals, the retro-cochlear auditory pathways are not seriously damaged and remain functionally intact even in the longer term after RT . Overall hearing outcomes after cochlear implantation for post- irradiated patients were found not to be worse than in patients who have had no prior RT to ear structures and related complications are rare, though so far, no study has compared the incidence of complications of post-cochlear implantation in irradiated versus non-irradiated temporal bones [195,196]. Moreover, cochlear implant surgery and cochlear implant activity have not been seen to have harmful effects on vestibular function and balance .
Current clinical standard includes proactive ototoxicity monitoring and standard aural rehabilitation with hearing aids for the vast majority of cases with efficacy of cochlear implants supported for profound SNHL after RT.
Conclusions
RT undeniably contributes to favorable disease control in the great majority of HNC patients and will remain one of the integral components of a multidisciplinary treatment approach for this disease. Even with recent improvements in RT planning and delivery with enhanced control over spatial distribution of radiation dose, some damage to normal tissues cannot be avoided . While the management of acute RT-induced toxicity is rather effective, the treatment options of late sequelae of RT are much more limited and of variable benefit which goes in line with the generally rather low quality level of available evidence (Table 4). Accordingly, there are no widely accepted clinical standards of management and treatment for these conditions, but only recommendations that could be used in clinic associated with variable degrees of benefit to the patients. The impact of late RT toxicity to the QOL in long-term HNC survivors cannot be overstated; indeed, they can also be life-threatening as suggested by long-term results of studies with mature follow-up [14,199]. Several new strategies are under development and/or extensive evaluation to improve the balance between efficacy and morbidity. Some of them refer to refinements in diagnostics (tumor identification) or RT techniques (e.g. guidelines for the delineation of target volumes and organ at risk, adaptive RT, biologically conformal RT/dose-painting, proton RT) ; other innovative approaches include a decrease in the treatment intensity for selected cases (e.g. HPV-associated oropharyngeal cancer), together with optimization in selection among different treatment options by using predictive tests in order to identify hypersensitive patients for RT-induced complications for a personalized strategy [201,202]. Enormous expansion in the field of innovative drug design offers a good prospect for the development of new preventive and therapeutic agents to counteract different RT-induced late sequelae.
Table 4.
Late sequeale after radiotherapy: recommended radiotherapy dose constrains, therapeutic options and their level of evidence (according to the American Society of Clinical Oncology Levels and Grades of Evidence)
| Toxicity | RT dose constraints (198) | Treatment options | Level of evidence |
:-- :-- |
| Xerostomia | Parotid gland: | Muscarinic agonist stimulation | Level II, Grade B |
| Dmean≤26 Gy or | Acupuncture | Level II, Grade B |
| V30<50% | Hyperbaric oxygen therapy | Level III, Grade B |
| Submandibular gland: | Gustatory/masticatory stimulants, mucosal lubricants and saliva substitutes | Level II, Grade B |
| Dmean<35 Gy |
|
|
| Dysphagia | Pharyngeal constricto muscle: | Exercise-based swallowing rehabilitation | Level II, Grade C |
| Dmean<50 Gy | Acupuncture | Level II, Grade B |
|
|
| Mandible osteoradionecrosis | Mandibular bone: | Hyperbaric oxygen therapy | Level II, grade B |
| 1 cc: Dmax 70–73.5 Gy | Medical management | Level III, Grade B |
|
|
| Trismus | Temporomandibular joint: | Exercise therapy | Level II; Grade B |
| 0.1 cc <70 Gy | Surgery | Level III, Grade B |
|
|
| Hearing loss | Inner ear: | Systemic steroids therapy | Level III, Grade B |
| Dmean<50 Gy | Cochlear implants | Level IV, Grade B |
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RT – Radiotherapy; D – Dose, VX – Volume that receives dose of X Gy.
Acknowledgments
Funding sources
This research did not receive any specific grant from funding agencies in the public, commercial, or not-for-profit sectors.
Footnotes
Conflict of interest statement
The authors have declared no conflicts of interest.
This article was written by members and invitees of the International Head and Neck Scientific Group (www.IHNSG.com).
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Learn About Family of Circles - Definitions, Formulas, Types & Solved Examples
Learn About Family of Circles - Definitions, Formulas, Types & Solved Examples
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The Family of Circles is a crucial concept in the world of coordinate geometry. It refers to a vast collection of circles, each sharing identical properties. These circles are grouped together in various ways to form what we call the family of circles. This article will guide you through the formulas related to the family of circles, supported by a few solved examples for better understanding.
The general equation of a circle is represented as x 2 + y 2 + 2gx + 2fy + c = 0. We will explore different ways of deriving the family of circles under specific conditions.
Let's denote S as x 2 + y 2 + 2gx + 2fy + c = 0
And S 1 as x 2 + y 2 + 2g 1 x + 2f 1 y + c 1 = 0
Similarly, S 2 can be represented as x 2 + y 2 + 2g 2 x + 2f 2 y + c 2 = 0
Let's denote L as lx + my + n = 0
Types of Circles in the Family of Circles
1. Family of circles with a fixed centre
This type of circle is represented by the equation (x-h)2 + (y-k)2 = r 2. Here, (h, k) is a fixed point, and r is a variable parameter. This equation represents a family of concentric circles. The fixation of the radius will result in a specific circle.
2. Family of circles passing through the intersection point of a line and a circle
For this type, the required equation for the family of circles passing through the intersection point of circle S = 0 and line L = 0 is given by S+λL = 0, where λ is a parameter.
3. Family of circles passing through the intersection point of two circles
The equation of the family of circles passing through the intersection point of two circles S 1 = 0 and S 2 = 0 is represented by S 1 + λS 2 = 0, where λ ≠ -1.
Remember, both circles should be in the given format of S 1 and S 2.
If λ = -1, the above equation becomes S 1 – S 2 = 0, which is the equation of the common chord if the circles intersect.
4. Family of circles passing through two given points
The equation of the family of circles passing through two given points P(x 1, y 1) and Q(x 2, y 2) is given by
(x−x 1)(x−x 2)+(y−y 1)(y−y 2)+λ∣∣∣∣x x 1 x 2 y y 1 y 2 1 1 1∣∣∣∣=0(x−x 1)(x−x 2)+(y−y 1)(y−y 2)+λ|x y 1 x 1 y 1 1 x 2 y 2 1|=0
5. Family of circles touching a line at a given point
The equation of a family of circles touching the line y-y 1 = m(x-x 1) at a point (x 1, y 1) on any finite m is given by
(x-x 1)2 + (y-y 1)2 + λ[(y-y 1) – m(x-x 1)] = 0
If m is infinite, the equation becomes (x – x 1)2 + (y-y 1)2 + λ(x-x 1) = 0
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Understanding Different Types of Circles
Orthogonal Circles
Two circles cutting each other at right angles are known as orthogonal circles. If the radii of two circles are r 1 and r 2 and d is the distance between their centres, then they are orthogonal if r 1 2 + r 2 2 = d 2.
Concentric Circles
Circles sharing the same centre but having different radii are termed as concentric circles. The equation of a circle with centre (-g, -f) and radius √(g 2 + f 2 -c) is given by x 2 + y 2 + 2gx + 2fy + c = 0. The equation of a concentric circle to the above circle is represented by x 2 + y 2 + 2gx + 2fy + c 1 = 0. Both the circles share the same centre (-g, -f), but they have different radii (c ≠ c 1).
Contact of Circles
Case 1: If the sum of the radii of two circles equals the distance between their centres, then the circles touch externally. The circles will satisfy the condition c 1 c 2 = r 1 + r 2.
Case 2: If the difference between the radii equals the distance between the centres, then the circles touch internally. The circles will satisfy the condition c 1 c 2 = r 1 – r 2.
Recommended Read
JEE Main Circles and System of Circles Previous Year Questions with Solutions
JEE Main Circle Previous Year Questions with Solutions
Solved Examples
Example 1: Let's consider the circles S 1 = x 2 + y 2 – 8x + 6y – 23 = 0 and S 2 = x 2 + y 2 – 2x – 5y + c = 0. Find the value of c if the circles S 1 and S 2 are orthogonal.
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Frequently Asked Questions
Give the general equation of a circle.
The general equation of a circle is given by x2 + y2 + 2gx + 2fy + c = 0.
What do you mean by orthogonal circles?
If two circles cut each other at right angles, then they are orthogonal circles.
What are concentric circles?
Concentric circles are circles having the same centre and different radius.
What is the condition for circles to touch externally?
The condition for the circles to touch externally is that the sum of the radii of two circles is equal to the distance between the centres.
What is the condition for circles to touch internally?
The condition for the circles to touch internally is that the difference between the radii and the distance between the centres are equal.
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13710 | https://chem.libretexts.org/Courses/University_of_British_Columbia/CHEM_100%3A_Foundations_of_Chemistry/08%3A_Quantities_in_Chemical_Reactions/8.3%3A_Mole-to-Mole_Conversions | Skip to main content
8.3: Mole-to-Mole Conversions
Last updated
: Sep 3, 2019
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8.2: Stoichiometry
8.4: Making Molecules: Mole to Mass (or vice versa) and Mass-to-Mass Conversions
Page ID
: 182683
( \newcommand{\kernel}{\mathrm{null}\,})
/;Learning Objectives
Use a balanced chemical equation to determine molar relationships between substances.
Previously, you learned to balance chemical equations by comparing the numbers of each type of atom in the reactants and products. The coefficients in front of the chemical formulas represent the numbers of molecules or formula units (depending on the type of substance). As follows, we will extend the meaning of the coefficients in a chemical equation.
Consider the simple chemical equation:
2H2+O2→2H2O
The convention for writing balanced chemical equations is to use the lowest whole-number ratio for the coefficients. However, the equation is balanced as long as the coefficients are in a 2:1:2 ratio. For example, this equation is also balanced if we write it as
4H2+2O2→4H2O
The ratio of the coefficients is 4:2:4, which reduces to 2:1:2. The equation is also balanced if we were to write it as
22H2+11O2→22H2O
because 22:11:22 also reduces to 2:1:2.
Suppose we want to use larger numbers. Consider the following coefficients:
12.044×1023H2+6.022×1023O2→12.044×1023H2O
These coefficients also have the ratio 2:1:2 (check it and see), so this equation is balanced. But 6.022 × 1023 is 1 mol, while 12.044 × 1023 is 2 mol (and the number is written that way to make this more obvious), so we can simplify this version of the equation by writing it as
2molH2+1molO2→2molH2O
We can leave out the word mol and not write the 1 coefficient (as is our habit), so the final form of the equation, still balanced, is
2H2+O2→2H2O
Now we interpret the coefficients as referring to molar amounts, not individual molecules. The lesson? Balanced chemical equations are balanced not only at the molecular level, but also in terms of molar amounts of reactants and products. Thus, we can read this reaction as “two moles of hydrogen react with one mole of oxygen to produce two moles of water.”
By the same token, the ratios we constructed to describe a molecular reaction can also be constructed in terms of moles rather than molecules. For the reaction in which hydrogen and oxygen combine to make water, for example, we can construct the following ratios:
2molH21molO2or1molO22molH2
2molH2O1molO2or1molO22molH2O
2molH22molH2Oor2molH2O2molH2
We can use these ratios to determine what amount of a substance, in moles, will react with or produce a given number of moles of a different substance. The study of the numerical relationships between the reactants and the products in balanced chemical reactions is called stoichiometry.
Example 8.3.1
How many moles of oxygen react with hydrogen to produce 27.6 mol of H2O?
Solution
Solutions to Example 8.3.1
| Steps for Problem Solving | How many moles of oxygen react with hydrogen to produce 27.6 mol of H2O? |
| Find a balanced equation that describes the reaction. | Unbalanced: H2 + O2 → H2O Balanced: 2H2 + O2 → 2H2O |
| Identify the "given" information and what the problem is asking you to "find." | Given: moles H2O Find: moles oxygen |
| List other known quantities. | 1 mol O2 = 2 mol H2O |
| Prepare a concept map and use the proper conversion factor. | |
| Cancel units and calculate. | 27.6molH2O×1molO22molH2O=13.8molO2 To produce 27.6 mol of H2O, 13.8 mol of O2 react. |
| Think about your result. | Since each mole of oxygen produces twice as many moles of water, it makes sense that the produced amount is greater than the reactant amount |
Example 8.3.2
How many moles of ammonia are produced if 4.20 moles of hydrogen are reacted with an excess of nitrogen?
Solution
Solutions to Example 8.3.2
| Steps for Problem Solving | How many moles of ammonia are produced if 4.20 moles of hydrogen are reacted with an excess of nitrogen? |
| Find a balanced equation that describes the reaction. | Unbalanced: N2 + H2 → NH3 Balanced: N2 + 3H2 → 2NH3 |
| Identify the "given" information and what the problem is asking you to "find." | Given: H2=4.20mol Find: mol of NH3 |
| List other known quantities. | 3 mol H2 = 2 mol NH3 |
| Prepare a concept map and use the proper conversion factor. | |
| Cancel units and calculate. | 4.20molH2×2molNH33molH2=2.80molNH3 The reaction of 4.20mol of hydrogen with excess nitrogen produces 2.80mol of ammonia. |
| Think about your result. | The result corresponds to the 3:2 ratio of hydrogen to ammonia from the balanced equation. |
Exercise 8.3.3
Given the following balanced chemical equation:
C5H12+8O2→5CO2+6H2O
, How many moles of H2O can be formed if 0.0652 mol of C5H12 were to react?
2. Balance the following unbalanced equation and determine how many moles of H2O are produced when 1.65 mol of NH3 react:
NH3+O2→N2+H2O
Answer a
: 0.391 mol H2O
Answer b
: 4NH3 + 3O2 → 2N2 + 6H2O; 2.48 mol H2O
Summary
The balanced chemical reaction can be used to determine molar relationships between substances.
8.2: Stoichiometry
8.4: Making Molecules: Mole to Mass (or vice versa) and Mass-to-Mass Conversions |
13711 | https://tlahuac.tecnm.mx/PDFs/automotrizcompetencias/6_SAF-1310.pdf | ©TecNM diciembre 2017 Página|1 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa 1. Datos Generales de la asignatura Nombre de la asignatura: Clave de la asignatura: SATCA1: Carrera: Diseño y Selección de Elementos de Máquinas SAF-1310 3-2-5 Ingeniería en Sistemas Automotrices 2. Presentación Caracterización de la asignatura Esta asignatura aporta al perfil del egresado en Sistemas Automotrices, la capacidad de diseñar y seleccionar diferentes elementos mecánicos de sujeción, ensamble y transmisión de potencia, para el desarrollo e implementación de sistemas automotrices. El contenido de esta signatura le permite formar competencias al egresado sobre el comportamiento de diversos elementos de máquina con cualquier tipo de cargas que se apliquen, calcular esfuerzos, determinar las características físicas de la pieza que soporte dicho esfuerzo, seleccionar diferentes elementos mecánicos, para ser capaz de optimizar el diseño de sistemas automotrices. Para cursar este programa se requieren competencias previas de las asignaturas Tecnología y Comportamiento de los Materiales, en lo que respecta a las características y propiedades de los diversos materiales, pruebas para la determinación de estás; Estática, diagrama de cuerpo libre, estática de la partícula, cuerpos rígidos, centros de gravedad y momentos de inercia, fuerzas en cables y armaduras; Mecánica de Materiales, esfuerzo y deformación, torsión, flexión y deformaciones; para la comprensión y aplicación de los temas que contempla el programa, Análisis y Síntesis de Mecanismos para la selección de elementos que harán la acción de junta, elementos de trasmisión de velocidad y/o potencia en un sistema automotriz. El contenido de la asignatura proporciona las bases para aplicar las herramientas computacionales de acuerdo a las tecnologías de vanguardia, en el diseño, simulación, operación y optimización de sistemas automotrices acordes a la demanda del sector industrial. Intención didáctica Se organiza el contenido en seis temas. En el tema uno se consideran los conceptos generales de la metodología y procedimientos para llevar a cabo un diseño mecánico de elementos y/o conjuntos aplicados al automóvil y a las maquinas que participan en el proceso de fabricación. 1 Sistema de Asignación y Transferencia de Créditos Académicos ©TecNM diciembre 2017 Página|2 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa El tema dos contempla conceptos generales de la concentración de esfuerzos y teorías de fallas, para el análisis y comprensión de los temas de diseño de sujetadores, diseño de engranes, selección de elementos mecánicos y el diseño de ejes de transmisión. El tema tres aborda los temas de elementos roscados, remaches, soldaduras como elementos de sujeción. El tema cuatro trata temas de los diversos tipos de engranes, referente a los tipos de carga y esfuerzos a que son sometidos los dientes, así como la aplicación de los engranes y trenes de engranajes en sistemas automotrices. El tema cinco hace referencia a elementos de transmisión como son: rodamientos, bandas, poleas, catarinas, cadenas, coples y reductores de velocidad, para poder analizar la aplicación en sistemas automotrices, utilizando normas nacionales e internacionales para asegurar que los elementos seleccionados cumplan con las normas de calidad del sector automotriz. Se sugiere el manejo de manuales y software proporcionado por los proveedores de dichos elementos para realizar la selección, montaje y mantenimiento de tales elementos de transmisión. El tema seis trata sobre el diseño de ejes de transmisión. En esta parte se analiza el procedimiento para el diseño de un eje en cuanto a carga estática y dinámica, verificando su velocidad crítica y las diversas aplicaciones en sistemas automotrices. 3. Participantes en el diseño y seguimiento curricular del programa Lugar y fecha de elaboración o revisión Participantes Evento Instituto Tecnológico Superior de Irapuato, del 13 al 16 de noviembre de 2012. Representantes de los Institutos Tecnológicos de: Apizaco, Celaya, Matamoros, Querétaro, Reynosa, Saltillo, San Juan del Río, San Luis Potosí, Tehuacán, Tepic, Tijuana, Tláhuac, Tláhuac II, Tlalnepantla, Superior de Lerdo, Superior de Libres, Superior del Sur de Guanajuato y Superior de Irapuato. Reunión Nacional de Diseño e Innovación Curricular para la Formación y Desarrollo de Competencias Profesionales de la carrera de Ingeniería en Sistemas Automotrices. ©TecNM diciembre 2017 Página|3 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa Desarrollo de Programas en Competencias Profesionales por los Institutos Tecnológicos del 19 de noviembre de 2012 al 1 de marzo de 2013. Academias de la carrera de Ingeniería en Sistemas Automotrices de los Institutos Tecnológicos de: Apizaco, Celaya, San Juan del Río, Tepic, Tláhuac, Superior de Irapuato y Superior de Libres. Elaboración del Programa de estudio propuesto en la Reunión Nacional de Diseño Curricular de la carrera de Ingeniería en Sistemas Automotrices. Instituto Tecnológico de Tláhuac, del 4 al 7 de marzo de 2013. Representantes de los Institutos Tecnológicos de: Apizaco, Celaya, Matamoros, Querétaro, Reynosa, Saltillo, San Juan del Río, San Luis Potosí, Tehuacán, Tepic, Tijuana, Tláhuac, Tláhuac II, Tlalnepantla, Superior de Lerdo, Superior de Libres y Superior de Irapuato. Reunión Nacional de Consolidación de los Programas en Competencias Profesionales de la carrera de Ingeniería en Sistemas Automotrices. Tecnológico Nacional de México, del 5 al 8 de diciembre de 2017. Representantes de los Institutos Tecnológicos de: Apizaco, Boca del Río, Superior de Abasolo, Superior de Lerdo, Superior de Irapuato, Superior de Libres y Superior del Oriente del Estado de Hidalgo. Reunión Nacional de Seguimiento Curricular de los Programas Educativos de Ingeniería en Animación Digital y Efectos Visuales, Ingeniería en Sistemas Automotrices y Licenciatura en Turismo. 4. Competencia(s) a desarrollar Competencia(s) específica(s) de la asignatura • Diseña diferentes sistemas de transmisión, flexible y de potencia, utilizados en maquinaria, equipo y sistemas automotrices, seleccionando los elementos adecuados para la aplicación requerida, así como el montaje y mantenimiento de tales elementos, para el funcionamiento óptimo de máquinas y sistemas. 5. Competencias previas • Calcula los esfuerzos y deformaciones que se presentan en los elementos mecánicos sujetos a cargas simples y combinadas, para interpretar los resultados. • Maneja instrumentos de medición, tolerancias geométricas y ajustes para la interpretación de planos de maquinaria y equipo, realizando conversiones entre sistemas de unidades y considerando las normas nacionales e internacionales. • Comprende las propiedades, procedimientos y el comportamiento de los diferentes materiales utilizados en ingeniería automotriz para controlar las características de los materiales con base en su microestructura, reconociendo los efectos en el medio ©TecNM diciembre 2017 Página|4 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa ambiente y las condiciones de operación sobre el rendimiento de los materiales. • Comprende los diversos procesos de manufactura para su utilización en la producción de elementos del sector automotriz. 6. Temario No. Temas Subtemas 1 Introducción al diseño en ingeniería 1.1 Introducción al diseño 1.2 Metodología de diseño 1.3 Ingeniería concurrente en el sector automotriz. 2 Teorías de falla 2.1 Modos de falla 2.2 Concentración de esfuerzos bajo carga estática y dinámica 2.3 Teorías de falla para materiales frágiles 2.4 Teorías de falla para materiales dúctiles 3 Diseño de sujetadores 3.1 Tornillos 3.2 Remaches 3.3 Soldaduras 3.4 Aplicación de sujetadores en sistemas automotrices 4 Diseño de engranes 4.1 Análisis de fuerzas en engranes rectos, helicoidales, cónicos y sinfín-corona. 4.2 Esfuerzos en dientes. 4.3 Normas y códigos de diseño. 4.4 Aplicaciones de engranes y trenes de engranajes en sistemas automotrices 5 Selección de elementos mecánicos 5.1 Tipos, aplicaciones, mantenimiento y selección de elementos mecánicos en sistemas automotrices 5.1.1 Rodamientos 5.1.2 Bandas y poleas 5.1.3 Cadenas y catarinas 5.1.4 Coples 5.1.5 Selección de reductores de velocidad 6 Diseño de ejes de transmisión 6.1 Procedimiento para el diseño de un eje 6.2 Diseño bajo carga estática y dinámica 6.3 Velocidad critica 6.4 Aplicaciones de los ejes de transmisión en sistemas automotrices ©TecNM diciembre 2017 Página|5 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa 7. Actividades de aprendizaje de los temas Tema 1. Introducción al Diseño en Ingeniería Competencias Actividades de aprendizaje Específica(s): • Interpreta las etapas de un diseño mecánico, para la correcta aplicación y ejecución del diseño. Genéricas: • Capacidad de análisis y síntesis • Habilidades de investigación. • Habilidad para buscar y analizar información proveniente de fuentes diversas. • Solución de problemas • Trabajo en equipo • Capacidad de aplicar los conocimientos en la práctica. • Revisar fuentes de información, interpretar y explicar los conceptos básicos y las generalidades que dan como metodología de un diseño mecánico. • Explicar la importancia que tiene cada una de las etapas de un diseño mecánico, y realizar un mapa conceptual. • Investigar en diversas fuentes de información, interpretar y explicar los conceptos básicos de la ingeniería concurrente y ejemplos de su impacto en el sector automotriz. • Analizar un caso de aplicación con la metodología para un diseño mecánico. Tema 2. Teorías de Falla Competencias Actividades de aprendizaje Específica(s): • Identifica y analiza las concentraciones de esfuerzo y modo de falla que se tiene de acuerdo a la solicitación de carga del material, geometría y el tipo de material utilizado, para optimizar diseños mecánicos con base a la geometría y tipo de material. Genéricas: • Capacidad de análisis y síntesis • Habilidades de investigación. • Habilidad para buscar y analizar información proveniente de fuentes diversas. • Solución de problemas • Trabajo en equipo • Capacidad de aplicar los conocimientos en la práctica. • Investigar en diversas fuentes de información los elementos causantes de concentración de esfuerzo en elementos mecánicos y realizar una exposición. • Determinar el estado de esfuerzo de una pieza cuando tiene secciones concentradoras de esfuerzos, bajo carga simple y combinada. • Investigar las teorías de fallas que son aplicables a los elementos mecánicos y realizar un mapa conceptual. • Comparar las teorías más comunes que existen para predecir la falla en materiales utilizados en elementos mecánicos y de acuerdo a los resultados obtenidos, discutir cual teoría es la más conservadora e idónea para el caso de análisis. ©TecNM diciembre 2017 Página|6 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa Tema 3. Diseño de Sujetadores Competencias Actividades de aprendizaje Específica(s): • Analiza, calcula y selecciona tornillos de potencia, remaches, y soldaduras como sujetadores sometidos a diferentes cargas y aplicaciones, para comprender su aplicación en sistemas automotrices Genéricas: • Capacidad de análisis y síntesis • Habilidades de investigación. • Habilidad para buscar y analizar información proveniente de fuentes diversas. • Solución de problemas • Trabajo en equipo • Capacidad de aplicar los conocimientos en la práctica • Explicar la clasificación y designación de los diferentes tipos de roscas. • Revisar planos técnicos de diferentes tipos de roscas. • Calcular el par de torsión, potencia, eficiencia en tornillos de potencia y auto bloqueo. • Calcular la carga, resistencia y parámetros de rigidez en pernos bajo carga estática y dinámica. • Seleccionar tornillos de potencia adecuados de acuerdo a los cálculos realizados. • Resolver problemas de diseño de juntas de empaquetadura. • Explicar la clasificación y designación de los diferentes tipos de remaches. • Revisar planos técnicos de diferentes tipos de remaches. • Calcular esfuerzos simples y combinados en uniones remachadas. • Desarrollar una Investigación de campo, donde se aprecie la importancia y aplicaciones de soldaduras. • Revisar planos técnicos de diferentes tipos de soldaduras. • Elaborar una tabla con las propiedades mínimas y los esfuerzos permisibles de las uniones soldadas con base a normas. • Revisar las fuentes de información correspondientes para establecer los criterios de diseño y decidir si la unión soldada es satisfactoria. • Resolver problemas donde se determine la resistencia de juntas soldadas bajo carga estática y bajo carga dinámica. ©TecNM diciembre 2017 Página|7 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa Tema 4. Diseño de Engranes Competencias Actividades de aprendizaje Específica(s): • Aplica la nomenclatura, normas y códigos de diseño utilizadas en los engranes y determina las fuerzas que afectan a los dientes de un engrane, enfocado al diseño de sistemas automotrices. Genéricas: • Capacidad de análisis y síntesis • Habilidades de investigación. • Habilidad para buscar y analizar información proveniente de fuentes diversas. • Solución de problemas • Trabajo en equipo • Capacidad de aplicar los conocimientos en la práctica. • Investigar las diferentes partes geométricas que componen a los engranes, revisando planos técnicos. • Analizar y determinar las fuerzas radiales y tangenciales, que actúan en los dientes de los engranes rectos, cónicos y helicoidales. • Determinar los esfuerzos admisibles en un diente de engrane, y la relación con el paso diametral, tamaño del diente y engrane, por medio de formulaciones como el criterio de Lewis. • Investigar en que normas y códigos se basa el diseño de engranes y realizar una discusión grupal. • Realizar una investigación de campo para conocer y analizar la aplicación de engranes y trenes de engranajes en sistemas automotrices. Tema 5. Selección de Elementos Mecánicos Competencias Actividades de aprendizaje Específica(s): • Analiza y selecciona elementos mecánicos involucrados en la transmisión de potencia como rodamientos, bandas, poleas, cadenas, catarinas, coples y reductores de velocidad, para, diseñar, optimizar sistemas automotrices. Genéricas: • Capacidad de análisis y síntesis • Habilidades de investigación. • Habilidad para buscar y analizar información proveniente de fuentes diversas. • Solución de problemas • Trabajo en equipo • Capacidad de aplicar los conocimientos en la práctica. • Investigar y realizar un resumen de los diferentes tipos de Rodamientos, clasificación, mantenimiento, montajes y aplicaciones. • Con base a manuales, software y catálogos de fabricantes, realizar la selección de diversos tipos de rodamientos sujetos a diferentes condiciones de trabajo, velocidad y carga, para determinar la vida útil, y capacidades de los rodamientos. • Investigar y realizar un resumen de los diferentes tipos de transmisiones flexibles de bandas y cadenas, clasificación, mantenimiento, montajes y aplicaciones. • Analizar y resolver problemas de sistemas que involucren el diseño y selección de la transmisión por bandas y poleas. • Analizar y resolver problemas de sistemas que involucren el diseño y ©TecNM diciembre 2017 Página|8 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa selección de la transmisión por cadenas y catarinas. • Investigar y realizar un resumen de los diferentes tipos de coples y Reductores de velocidad, clasificación, mantenimiento, montajes y aplicaciones. • Analizar y resolver problemas de sistemas que involucren la transmisión de potencia mediante la selección de coples. • Analizar y resolver problemas de sistemas que involucren la transmisión de potencia mediante la selección de reductores de velocidad. • Realizar una investigación de campo, sobre las principales consecuencias y causas de falla provocadas por una mala selección de elementos mecánicos de transmisión de potencia. Tema 6. Diseño de ejes de transmisión Competencias Actividades de aprendizaje Específica(s): • Aplica los principios básicos para el diseño de ejes sujetos a cargas estáticas y cíclicas, así como determina la primera velocidad crítica de un eje, para diseñar, mantener y optimizar ejes de transmisión en el sector automotriz. Genéricas: • Capacidad de análisis y síntesis • Habilidades de investigación. • Habilidad para buscar y analizar información proveniente de fuentes diversas. • Solución de problemas • Trabajo en equipo • Capacidad de aplicar los conocimientos en la práctica. • Investigar y explicar la importancia y aplicaciones del diseño de ejes, así como los procedimientos y análisis que esto involucra. • Analizar planos técnicos de ejes de transmisión para discutir el ensamble y elementos mecánicos involucrados. • Dibujar un eje donde se puedan observar las cargas y esfuerzos a los que se encuentra sujeto y los diagramas de par torsional, momento horizontal y momento vertical. • Resolver problemas de diseño de ejes sujetos a carga estática, bajo carga axial, momento flexionante y torsión, aplicando las teorías de la Energía de Distorsión y del Esfuerzo Cortante Máximo. • Resolver problemas de diseño de ejes sujetos a carga cíclica, tanto de momento flexionante alternante, torsión continua, torsión alternante, aplicando las teorías de: Energía de Distorsión y Esfuerzo Cortante Máximo, para ©TecNM diciembre 2017 Página|9 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa materiales dúctiles; y el criterio de Esfuerzo Normal Máximo para materiales frágiles. • Explicar la importancia del cálculo de la primera velocidad crítica de un eje, sus causas, análisis, medición, prevención y efectos en el diseño y funcionamiento del eje. • Realizar una investigación de campo para conocer y analizar la aplicación de los ejes de transmisión en sistemas automotrices. 8. Práctica(s) • Analizar mediante prototipos didácticos concentraciones de esfuerzos en diferentes elementos mecánicos aplicando métodos experimentales, tales como: extensometría y fotoelasticidad. • Analizar y resolver problemas de concentraciones de esfuerzos en elementos mecánicos aplicando software de elemento finito. • Realizar pruebas en máquina universal con cargas constantes y repetitivas para analizar la resistencia y esfuerzos en uniones atornilladas, remachadas y soldadas en prototipos didácticos. • Seleccionar diversos tipos de rodamientos mediante software de fabricantes. • Realizar el montaje y mantenimiento de sistemas de transmisión de banda y polea, cadena y catarina, coples y reductores de velocidad. • Calcular y determinar experimentalmente la primera velocidad crítica de un eje, explicar sus causas, medición, prevención y efectos en el diseño y funcionamiento del eje. • Realizar visitas industriales con la finalidad de observar las aplicaciones del diseño y selección de elementos mecánicos en el sector automotriz. Con la presentación del correspondiente. 9. Proyecto de asignatura El objetivo del proyecto que planteé el docente que imparta esta asignatura, es demostrar el desarrollo y alcance de la(s) competencia(s) de la asignatura, considerando las siguientes fases: • Fundamentación: marco referencial (teórico, conceptual, contextual, legal) en el cual se fundamenta el proyecto de acuerdo con un diagnóstico realizado, mismo que permite a los estudiantes lograr la comprensión de la realidad o situación objeto de estudio para definir un proceso de intervención o hacer el diseño de un modelo. • Planeación: con base en el diagnóstico en esta fase se realiza el diseño del proyecto por parte de los estudiantes con asesoría del docente; implica planificar un proceso: de intervención empresarial, social o comunitario, el diseño de un modelo, entre otros, según el tipo de proyecto, las actividades a realizar los recursos requeridos y el cronograma de trabajo. ©TecNM diciembre 2017 Página|10 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa • Ejecución: consiste en el desarrollo de la planeación del proyecto realizada por parte de los estudiantes con asesoría del docente, es decir en la intervención (social, empresarial), o construcción del modelo propuesto según el tipo de proyecto, es la fase de mayor duración que implica el desempeño de las competencias genéricas y especificas a desarrollar. • Evaluación: es la fase final que aplica un juicio de valor en el contexto laboral-profesión, social e investigativo, ésta se debe realizar a través del reconocimiento de logros y aspectos a mejorar se estará promoviendo el concepto de “evaluación para la mejora continua”, la meta cognición, el desarrollo del pensamiento crítico y reflexivo en los estudiantes. 10. Evaluación por competencias Instrumentos y herramientas sugeridas para evaluar las actividades de aprendizaje: La evaluación se realiza con el propósito de evidenciar el desarrollo de las competencias específicas y genéricas de manera integral, creando las condiciones en distintos espacios de aprendizaje y desempeño profesional. En el contexto de la evaluación por competencias, dentro de las evidencias de desempeño, se sugieren las siguientes: • Mapas • Diagramas • Tabla comparativa • Ensayos • Evaluación • Cuadro sinóptico • Foros de discusión • Videos • Reportes • Bitácora • Resumen • Presentaciones Y los instrumentos de evaluación del desarrollo de competencias específicas y genéricas, pueden ser: • Guía de observación • Matriz de valoración • Lista de cotejo • Guía de proyectos • Rúbricas ©TecNM diciembre 2017 Página|11 TECNOLÓGICO NACIONAL DE MÉXICO Secretaría Académica, de Investigación e Innovación Dirección de Docencia e Innovación Educativa 11. Fuentes de información 1. Hamrock Bernard J., Jacobson Bo., Schmid Steven.(2000)” Elementos de máquinas”. Editorial Mc Graw Hill. 2. Budynas Richard G. (2012). “Diseño en Ing. Mecánica de Shigley”. Editorial Mc Graw Hill. 3. L Mott. R. (2006) “Diseño de Elementos de Máquinas”. Pearson Educación. 4. Spotts M. F., T. E. Shoup. (2001) “Elementos de máquinas”. Editorial Mc Graw Hill. 5. Collins Jack A. Busby Henry R., H. Staab George (2009) 1. “Mechanical Design of Machine Elements and Machines” 6. Spots Merhyle R. E. Shoup Terry. E. Hornberger L.(2011) “Design of Machine Elements” Hardcover. 7. Moring Faires.V. (2004) “Diseño de elementos de máquinas”. Editorial Uteha. 8. Juvinall, R.C. (2012). “Diseño de elementos de máquinas”. Limusa, 9. Norton, R.L. (1998) “Machine design”. New jersey, Editorial Prentice Hall. 10. SKF. (2010). Catálogo general de rodamientos” 11. TIMKEN. (2009) “Catálogo general de rodamientos” 12. DODGE. (2008). “Manual de selección para bandas” 13. GATES. (2009). “Manual de selección para bandas” 14. GATES No. 14955 – A, 8/99.”Manual de selección para bandas “V” y servicio pesado” |
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How to Calculate the Gravitational Force between Two Everyday Objects
High School Physics Skills Practice
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0:12 Example 1 -…
3:15 Example 2 -…
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Daniel Jibson, Kirsten Wordeman
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Daniel JibsonDaniel has taught physics and engineering since 2011. He has a BS in physics-astronomy from Brigham Young University and an MA in science education from Boston University. He currently holds a science teaching license for grades 8-12.
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Kirsten WordemanKirsten has taught high school biology, chemistry, physics, and genetics/biotechnology for three years. She has a Bachelor's in Biochemistry from The University of Mount Union and a Master's in Biochemistry from The Ohio State University. She holds teaching certificates in biology and chemistry.
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Example Solutions Practice Questions
Steps for Calculating the Gravitational Force between Two Everyday Objects
Step 1: Make a list including the mass of each object and the distance separating them.
Step 2: Use the formula for universal gravitation to calculate the force between the objects.
Vocabulary and Equations for Calculating the Gravitational Force between Two Everyday Objects
Newton's Law of Universal Gravitation: Isaac Newton determined that any two objects will exert a gravitational force upon each other proportional to the mass of both objects ({eq}m_1 {/eq} and {eq}m_2 {/eq}), and inversely proportional to the square of the distance of separation {eq}(r) {/eq}. The equation is:
$$F_g = \dfrac{Gm_1m_2}{r^2} $$
In this equation, {eq}G {/eq} is the universal gravitational constant: {eq}G = 6.67 \times 10^{-11}\ \rm{N\cdot m^2/kg^2} {/eq}.
The following two problems demonstrate how to calculate the gravitational force between two objects.
Example Problem 1 - Calculating the Gravitational Force between Two Everyday Objects
Calculate the gravitational force between two people, 62 kg and 75 kg, respectively, separated by 2.5 meters.
Step 1: Make a list including the mass of each object and the distance separating them.
The masses of both people and the distance separating them are given in the problem:
{eq}m_1 = 62\ \rm{kg} {/eq}
{eq}m_2 = 75\ \rm{kg} {/eq}
{eq}r = 2.5\ \rm{m} {/eq}
Step 2: Use the formula for universal gravitation to calculate the force between the objects.
We will apply Newton's law of universal gravitation to calculate the force between them:
$$F_g = \dfrac{Gm_1m_2}{r^2} = \dfrac{(6.67 \times 10^{-11}\ \rm{N\cdot m^2/kg^2})( 62\ \rm{kg})(75\ \rm{kg})}{(2.5\ \rm{m})^2} = 5.0\times 10^{-8}\ \rm{N} $$
The gravitational force between these two people is 0.000000050 Newtons.
Example Problem 2 - Calculating the Gravitational Force between Two Everyday Objects
A penny having a mass of 2.5 grams and a glass of water with a mass of 0.58 kg are 23 cm apart. What is the gravitational force between them?
Step 1: Make a list including the mass of each object and the distance separating them.
We have the following information:
{eq}m_1 = 2.5\ \rm{g} \times \dfrac{1\ \rm{kg}}{1000\ \rm{g}} = 0.0025\ \rm{kg} {/eq}
{eq}m_2 = 0.58\ \rm{kg} {/eq}
{eq}r = 23\ \rm{cm} \times \dfrac{1\ \rm{m}}{100\ \rm{cm}} = 0.23\ \rm{m} {/eq}
Notice that the units of grams and centimeters have been converted to kilograms and meters. This is necessary because of the units built into the universal gravitational constant.
Step 2: Use the formula for universal gravitation to calculate the force between the objects.
We will use the equation to calculate the gravitational force between the two objects.
$$F_g = \dfrac{Gm_1m_2}{r^2} = \dfrac{(6.67 \times 10^{-11}\ \rm{N\cdot m^2/kg^2})(0.0025\ \rm{kg})(0.58\ \rm{kg})}{(0.23\ \rm{m})^2} = 1.8\times 10^{-12}\ \rm{N} $$
The gravitational force between the penny and the glass of water is 0.0000000000018 N.
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Table of Contents
Steps for Calculating the Gravitational Force between Two Everyday Objects
Vocabulary and Equations for Calculating the Gravitational Force between Two Everyday Objects
Example Problem 1 - Calculating the Gravitational Force between Two Everyday Objects
Example Problem 2 - Calculating the Gravitational Force between Two Everyday Objects
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Calculating the Gravitational Force between Two Everyday Objects
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13714 | https://es.wikipedia.org/wiki/Diferenciaci%C3%B3n_de_producto | Diferenciación de producto - Wikipedia, la enciclopedia libre
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Inicio
1 Razón fundamental
2 Historia
3 Diferenciación vertical del producto
4 Diferenciación horizontal del producto
5 Productos sustitutos y diferenciación de productos
6 Interacción entre la diferenciación horizontal y vertical: Una aplicación a la banca
7 Véase también
8 Referencias
9 Enlaces externos
Diferenciación de producto
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En economía y marketing, la diferenciación de productos (o simplemente la diferenciación) es el proceso de distinguir un producto o servicio de otros, para hacerlo más atractivo para un mercado objetivo en particular. Esto implica diferenciarlo de los productos de la competencia y de los propios productos de la empresa. El concepto fue propuesto por Edward Chamberlin en su 1933 La Teoría de la Competencia Monopolística.
Razón fundamental
[editar]
Las firmas tienen diferentes dotaciones de recursos que les permiten construir ventajas competitivas específicas sobre los competidores. Las dotaciones de recursos permiten que las empresas sean diferentes, lo que reduce la competencia y hace posible alcanzar nuevos segmentos del mercado. Por lo tanto, la diferenciación es el proceso de distinguir las diferencias de un producto u oferta de otros, para hacerlo más atractivo para un mercado objetivo particular.
Aunque la investigación en un nicho de mercado puede dar como resultado el cambio de un producto para mejorar la diferenciación, los cambios en sí mismos no son diferenciación. La comercialización o la diferenciación del producto es el proceso de describir las diferencias entre los productos o los servicios, o la lista resultante de diferencias. Esto se hace para demostrar los aspectos únicos del producto de una empresa y crear un sentido de valor. Los libros de texto de la comercialización son firmes en el punto que cualquier diferenciación debe ser valorada por los compradores (un intento de diferenciación que no se percibe no cuenta). El término propuesta única de venta se refiere a la publicidad para comunicar la diferenciación de un producto.
En economía, la diferenciación exitosa de productos conduce a una ventaja competitiva y es incompatible con las condiciones para una competencia perfecta, que incluyen el requisito de que los productos de las firmas competidoras sean sustitutos perfectos. Existen tres tipos de diferenciación de producto:
Simple: basado en una variedad de características
Horizontal: basado en una única característica, pero los consumidores no tienen clara la calidad
Vertical: basado en una característica única y los consumidores son claros en cuanto a su calidad
Las diferencias de marca generalmente son menores; pueden ser simplemente una diferencia en el empaquetado o un tema publicitario. El producto físico no necesita cambiar, pero puede. La diferenciación se debe a que los compradores perciben una diferencia; por lo tanto, las causas de la diferenciación pueden ser aspectos funcionales del producto o servicio, cómo se distribuyen y comercializan, o quién la compra. Las principales fuentes de diferenciación de productos son las siguientes:
Diferencias de calidad que suelen ir acompañadas de diferencias de precio.
Diferencias en características funcionales o diseño.
Ignorancia de los compradores con respecto a las características y cualidades esenciales de los bienes que están comprando.
Actividades de promoción de ventas de vendedores y, en particular, publicidad.
Diferencias de disponibilidad (por ejemplo, tiempo y ubicación).
El objetivo de la diferenciación es desarrollar una posición que los clientes potenciales vean como única. El término se usa frecuentemente cuando se trata de modelos de negocios freemium, en los cuales las empresas comercializan una versión gratuita y pagada de un producto dado. Dado que se dirigen al mismo grupo de clientes, es imperativo que las versiones gratuitas y pagadas sean efectivamente diferenciadas.
La diferenciación afecta principalmente al rendimiento mediante la reducción de la franqueza de la competencia: a medida que el producto se vuelve más diferente, la categorización se hace más difícil y por lo tanto, dibuja menos comparaciones con su competencia. Una estrategia de diferenciación de producto exitosa moverá su producto de la competencia basada principalmente en el precio a la competencia en los factores que no son de precio (tales como características del producto, estrategia de distribución, o variables de promoción).
La mayoría de la gente diría que la implicación de la diferenciación es la posibilidad de cobrar un precio premium; sin embargo, se trata de una simplificación bruta. Si los clientes valoran la oferta de la empresa, serán menos sensibles a los aspectos de las ofertas competidoras; el precio puede no ser uno de estos aspectos. La diferenciación hace que los clientes en un segmento dado tengan una sensibilidad más baja a otras características (sin precio) del producto.
Historia
[editar]
La obra seminal de Edward Chamberlin (1933) sobre la competencia monopolística mencionó la teoría de la diferenciación que dice que para los productos disponibles dentro de la misma industria, los clientes pueden tener preferencias diferentes. Sin embargo, una estrategia genérica de diferenciación que fue popularizada por Michael Porter (1980) que es cualquier producto (tangible o intangible) percibido como "ser único" por al menos un conjunto de clientes. Por lo tanto, depende de su percepción el alcance de la diferenciación de los productos. Incluso hasta 1999, las consecuencias de estos conceptos no fueron bien entendidas. De hecho, Miller (1986) propuso la comercialización y la innovación como dos estrategias de diferenciación, que fue apoyado por algunos eruditos como Lee y Miller (1999). Mintzberf (1988) propuso categorías más específicas pero amplias: calidad, diseño, apoyo, imagen, precio y productos no diferenciados, que obtuvo el apoyo de Kotha y Vadlamani (1995). Sin embargo, la literatura de IO (Ethiraj & Zhu, 2008; Makadok, 2010, 2011) hizo un análisis más profundo de la teoría y exploró una clara distinción entre el uso amplio de la diferenciación vertical y horizontal.
Diferenciación vertical del producto
[editar]
Si los productos A y B tienen el mismo precio para el consumidor, la cuota de mercado para cada uno será positiva, según el modelo de Hotelling. La principal teoría de esto es que todos los consumidores prefieren el producto de mayor calidad si se ofrecen dos productos distintos al mismo precio. Un producto puede diferir en muchos atributos verticales, como su velocidad de funcionamiento. Lo que realmente importa es la relación entre la disposición de los consumidores a pagar por mejoras en la calidad y el aumento en el costo por unidad que conlleva tales mejoras. Por lo tanto, la diferencia percibida en la calidad es diferente con diferentes consumidores, por lo que es objetivo. Un producto verde puede estar teniendo un efecto negativo más bajo o cero en el ambiente, sin embargo, puede resultar ser inferior que otros productos en otros aspectos. Por lo tanto, también depende de la forma en que se anuncia y la presión social que un consumidor potencial está viviendo. Incluso una diferenciación vertical puede ser un factor decisivo en la compra.
Diferenciación horizontal del producto
[editar]
Cuando los productos no pueden ser ordenados de una manera objetiva y son diferentes en una o todas sus características, entonces hay diferenciación horizontal. Por ejemplo, diferentes versiones de color del mismo iPhone o MacBook. Un helado de limón no es superior a un helado de chocolate, se basa completamente en la preferencia del usuario. Un restaurante puede poner precio a todos sus postres con el mismo precio y permite al consumidor elegir libremente sus preferencias ya que todas las alternativas cuestan lo mismo.
Productos sustitutos y diferenciación de productos
[editar]
Según una investigación realizada mediante la combinación de matemáticas y economía, las decisiones de fijación de precios dependen de la substitución entre los productos depende del grado de diferenciación de los productos de las firmas. Ninguna empresa puede cobrar un precio más alto si los productos son buenos sustitutos y viceversa. Cuanto menor sea el precio de equilibrio no cooperativo, menor será la diferenciación. Por esta razón, las empresas podrían elevar conjuntamente los precios por encima del equilibrio o el nivel competitivo mediante la coordinación entre ellos mismos. Tienen un acuerdo de colusión verbal o escrito entre ellos. Las empresas que operan en un mercado de baja diferenciación de productos podrían no coordinarse con otras, lo que aumenta el incentivo para engañar el acuerdo de colusión. En el contraste, incluso bajando ligeramente los precios, una empresa puede capturar una gran fracción del mercado y obtener ganancias a corto plazo con productos altamente sustituibles.
Interacción entre la diferenciación horizontal y vertical: Una aplicación a la banca
[editar]
Durante los años 1990, las medidas adoptadas por el gobierno sobre la desreglamentación y la integración europea persuadieron a los bancos a competir por depósitos en muchos factores como las tasas de depósito, la accesibilidad y la calidad de los servicios financieros.
En este ejemplo usando el modelo de Hotelling, una característica es de variedad (ubicación) y una característica de calidad (acceso remoto). Acceso remoto mediante servicios bancarios a través de servicios postales y telefónicos como la organización de facilidades de pago y obtención de información de cuentas. En este modelo, los bancos no se pueden diferenciar verticalmente sin afectar negativamente la diferenciación horizontal entre ellos.
La diferenciación horizontal ocurre con la ubicación de la sucursal del banco. La diferenciación vertical, en este ejemplo, se produce siempre que un banco ofrezca acceso remoto y el otro no. Con el acceso alejado, puede estimular una interacción negativa entre la tarifa del transporte y el gusto para la calidad: los clientes que tienen un gusto más alto para el acceso alejado enfrentan una tarifa más baja del transporte.
Un depositante con un alto sabor (bajo) para el acceso remoto tiene bajos (altos) costos de transporte lineal. Diferentes equilibrios emergen como resultado de dos efectos. Por un lado, la introducción del acceso remoto roba a los depositantes de su competidor porque la especificación del producto se vuelve más atractiva (efecto directo). Por otra parte, los bancos se convierten en sustitutos más cercanos (efecto indirecto). Primero, los bancos se convierten en sustitutos más cercanos a medida que disminuye el impacto de los costos de transporte lineal. En segundo lugar, la competencia de tasa de depósito se ve afectada por el tamaño de la diferencia de calidad. Estos dos efectos, "robar" a los depositantes versus "sustituibilidad" entre los bancos, determinan el equilibrio. Para los valores bajos y altos de la diferencia de la calidad del cociente a la tarifa del transporte, solamente un banco ofrece el acceso alejado (especialización). Los valores intermedios (muy bajos) de la diferencia de calidad de la relación con los costos de transporte rinden acceso universal (no) a distancia.
Este concurso es un juego de dos factores: uno es de ofrecimiento de acceso remoto y el otro es de tasas de depósito. Hipotéticamente, habrá dos escenarios consecuentes si sólo un banco ofrece acceso remoto. En primer lugar, el banco obtiene una cuota de mercado positiva para todos los tipos de acceso remoto, dando lugar a un dominio horizontal. Esto ocurre cuando el costo de transporte prevalece sobre la calidad del servicio, la tasa de depósito y el tiempo. En segundo lugar, el dominio vertical aparece en la imagen cuando el banco que no está ofreciendo acceso remoto obtiene todo el mercado para los depositantes que tienen la menor preferencia por el acceso remoto. Es cuando el servicio de calidad, la tasa de depósito y el tiempo prevalecen sobre el costo del transporte.
Véase también
[editar]
Crippleware
Competencia sin precio
Marketing
Personalización en masa
Configurador
Segmentación de mercado
Gestión de producto
Marca
País de origen
Plan de marketing
Freemium
Posicionamiento
Referencias
[editar]
↑Chamberlin, Edward (1962). The Theory of Monopolistic Competition: A Re-orientation of the Theory of Value(en inglés). Harvard University Press. ISBN978-0674881259.
↑Barney, J (Marzo de 1991). «Firm resources and sustained competitive advantage»(en inglés). Journal of Management. pp.17 (1): 99-120. Consultado el 12 de abril de 2018.
↑ Saltar a: abKotler, Keller, Philip (2006). Marketing Management(en inglés). Prentice Hall New Jersey. ISBN9780130156846.
↑Reeves, Rosser (1961). The Reality of Advertising(en inglés). ISBN978-0982694145.
↑Pepall, Lynne; Daniel J. Richards; George Norman (2005). Industrial Organization, Contemporary Theory and Practice(en inglés). Natorp Boulevard, Mason, Ohio 45040: Thomson South-Western. p.133. ISBN0-324-22474-5.
↑Sharp, Byron; Dawes, John (2001). «What is Differentiation and How Does it Work?»(en inglés). Journal of Marketing Management. pp.17, 739-59.
↑«Being better vs. being different: Differentiation, competition, and pricing strategies in the Spanish hotel industry». Tourism Management(en inglés)34: 71-79. 1 de febrero de 2013. ISSN0261-5177. doi:10.1016/j.tourman.2012.03.014. Consultado el 10 de abril de 2018.
↑Sutton, John (1986). «Vertical Product Differentiation: Some Basic Themes». The American Economic Review76 (2): 393-398. Consultado el 10 de abril de 2018.
↑«Product differentiation - a key concept in Economics and Management». www.economicswebinstitute.org. Consultado el 10 de abril de 2018.
↑«Product differentiation - a key concept in Economics and Management». www.economicswebinstitute.org. Consultado el 10 de abril de 2018.
↑«Analysis of pricing decision for substitutable and complementary products with a common retailer». Pacific Science Review A: Natural Science and Engineering(en inglés)18 (3): 190-202. 1 de noviembre de 2016. ISSN2405-8823. doi:10.1016/j.psra.2016.09.012. Consultado el 11 de abril de 2018.
↑ Saltar a: abcdeDegryse, Hans (1996). «On the Interaction Between Vertical and Horizontal Product Differentiation: An Application to Banking». The Journal of Industrial Economics44 (2): 169-186. doi:10.2307/2950644. Consultado el 11 de abril de 2018.
Enlaces externos
[editar]
Primavera de 1997 - Discurso de Jonathan B. Baker Director, Bureau of Economics Federal Trade Commission sobre Diferenciación el Producto
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13715 | https://en.wikipedia.org/wiki/Sodium_dithionite | Jump to content
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From Wikipedia, the free encyclopedia
"Sodium hydrosulfite" redirects here. For sodium hydrogen sulfite (NaHSO3), sometimes erroneously referred to as sodium hydrosulfite, see Sodium bisulfite.
Not to be confused with Sodium hydrosulfide or Sodium dithionate.
"D-Ox" redirects here. For other uses, see dox (disambiguation).
Sodium dithionite
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| |
| Names |
| Other names D-Ox, Hydrolin, Reductonesodium hydrosulfite, sodium sulfoxylate, SulfoxylateVatrolite, Virtex LHydrosulfit, PrayonBlankit, Albite A, KoniteZepar, Burmol, Arostit |
| Identifiers |
| CAS Number | 7775-14-6Y |
| 3D model (JSmol) | Interactive image |
| ChEBI | CHEBI:66870N |
| ChemSpider | 22897 |
| ECHA InfoCard | 100.028.991 |
| EC Number | 231-890-0 |
| PubChem CID | 24489 |
| RTECS number | JP2100000 |
| UNII | 2K5B8F6ES1Y |
| UN number | 1384 |
| CompTox Dashboard (EPA) | DTXSID9029697 |
| InChI InChI=1S/2Na.H2O4S2/c;;1-5(2)6(3)4/h;;(H,1,2)(H,3,4)/q2+1;/p-2 Key: JVBXVOWTABLYPX-UHFFFAOYSA-L |
| SMILES [O-]S(=O)S(=O)[O-].[Na+].[Na+] |
| Properties |
| Chemical formula | Na2S2O4 |
| Molar mass | 174.107 g/mol (anhydrous) 210.146 g/mol (dihydrate) |
| Appearance | white to grayish crystalline powder light-lemon colored flakes |
| Odor | faint sulfur odor |
| Density | 2.38 g/cm3 (anhydrous) 1.58 g/cm3 (dihydrate) |
| Melting point | 52 °C (126 °F; 325 K) |
| Boiling point | Decomposes |
| Solubility in water | 18.2 g/100 mL (anhydrous, 20 °C) 21.9 g/100 mL (Dihydrate, 20 °C) |
| Solubility | slightly soluble in alcohol |
| Hazards |
| GHS labelling: |
| Pictograms | |
| Signal word | Danger |
| Hazard statements | H251, H302 |
| Precautionary statements | P235+P410, P264, P270, P280, P301+P312, P330, P407, P413, P420, P501 |
| NFPA 704 (fire diamond) | 2 3 1 |
| Flash point | 100 °C (212 °F; 373 K) |
| Autoignitiontemperature | 200 °C (392 °F; 473 K) |
| Related compounds |
| Other anions | Sodium sulfiteSodium sulfate |
| Related compounds | Sodium thiosulfateSodium bisulfiteSodium metabisulfiteSodium bisulfate |
| Except where otherwise noted, data are given for materials in their standard state (at 25 °C [77 °F], 100 kPa). N verify (what is YN ?) Infobox references |
Chemical compound
Sodium dithionite (also known as sodium hydrosulfite) is a white crystalline powder with a sulfurous odor. Although it is stable in dry air, it decomposes in hot water and in acid solutions.
Structure
[edit]
The structure has been examined by Raman spectroscopy and X-ray crystallography. The dithionite dianion has C2 symmetry, with almost eclipsed with a 16° O-S-S-O torsional angle. In the dihydrated form (Na2S2O4·2H2O), the dithionite anion has gauche 56° O-S-S-O torsional angle.
A weak S-S bond is indicated by the S-S distance of 239 pm, which is elongated by ca. 30 pm relative to a typical S-S bond. Because this bond is fragile, the dithionite anion dissociates in solution into the [SO2]− radicals, as has been confirmed by EPR spectroscopy. It is also observed that 35S undergoes rapid exchange between S2O42− and SO2 in neutral or acidic solution, consistent with the weak S-S bond in the anion.
Preparation
[edit]
Sodium dithionite is produced industrially by reduction of sulfur dioxide. Approximately 300,000 tons were produced in 1990. The route using zinc powder is a two-step process:
: 2 SO2 + Zn → ZnS2O4
: ZnS2O4 + 2 NaOH → Na2S2O4 + Zn(OH)2
The sodium borohydride method obeys the following stoichiometry:
: NaBH4 + 8 NaOH + 8 SO2 → 4 Na2S2O4 + NaBO2 + 6 H2O
Each equivalent of H− reduces two equivalents of sulfur dioxide. Formate has also been used as the reductant.
Properties and reactions
[edit]
Hydrolysis
[edit]
Sodium dithionite is stable when dry, but aqueous solutions deteriorate due to the following reaction:
: 2 S2O42− + H2O → S2O32− + 2 HSO3−
This behavior is consistent with the instability of dithionous acid. Thus, solutions of sodium dithionite cannot be stored for a long period of time.
Anhydrous sodium dithionite decomposes to sodium sulfate and sulfur dioxide above 90 °C in the air. In absence of air, it decomposes quickly above 150 °C to sodium sulfite, sodium thiosulfate, sulfur dioxide and trace amount of sulfur.
Redox reactions
[edit]
Sodium dithionite is a reducing agent. At pH 7, the potential is -0.66 V compared to the normal hydrogen electrode. Redox occurs with formation of bisulfite:
: S2O42- + 2 H2O → 2 HSO3− + 2 e− + 2 H+
Sodium dithionite reacts with oxygen:
: Na2S2O4 + O2 + H2O → NaHSO4 + NaHSO3
These reactions exhibit complex pH-dependent equilibria involving bisulfite, thiosulfate, and sulfur dioxide.
With organic carbonyls
[edit]
In the presence of aldehydes, sodium dithionite reacts either to form α-hydroxy-sulfinates at room temperature or to reduce the aldehyde to the corresponding alcohol above a temperature of 85 °C. Some ketones are also reduced under similar conditions.
Uses
[edit]
Industry
[edit]
Sodium dithionite is used as a water-soluble reducing agent in some industrial dyeing processes. In the case of sulfur dyes and vat dyes, an otherwise water-insoluble dye can be reduced into its water-soluble alkali metal leuco salt. Indigo dye is sometimes processed in this way.
Domestic and hobby uses
[edit]
Sodium dithionite can also be used for water treatment, aquarium water conditioners, gas purification, cleaning, and stripping. In addition to the textile industry, this compound is used in industries concerned with leather, foods, polymers, photography, and many others, often as a decolourising agent. It is even used domestically as a decolouring agent for white laundry, when it has been accidentally stained by way of a dyed item slipping into the high temperature washing cycle. It is usually available in 5 gram sachets termed hydrosulphite after the antiquated name of the salt.
It is the active ingredient in "Iron Out Rust Stain Remover", a commercial rust product.
Laboratory
[edit]
Sodium dithionite is often used in physiology experiments as a means of lowering solutions' redox potential (Eo' -0.66 V vs SHE at pH 7). Potassium ferricyanide is usually used as an oxidizing chemical in such experiments (Eo' ~ .436 V at pH 7). In addition, sodium dithionite is often used in soil chemistry experiments to determine the amount of iron that is not incorporated in primary silicate minerals. Hence, iron extracted by sodium dithionite is also referred to as "free iron."
Aqueous solutions of sodium dithionite were once used to produce 'Fieser's solution' for the removal of oxygen from a gas stream. Pyrithione can be prepared in a two-step synthesis from 2-bromopyridine by oxidation to the N-oxide with a suitable peracid followed by substitution using sodium dithionite to introduce the thiol functional group.
Photography
[edit]
It is used in Kodak fogging developer, FD-70. This is used in the second step in processing black and white positive images, for making slides. It is part of the Kodak Direct Positive Film Developing Outfit.
Safety
[edit]
The wide use of sodium dithionite is attributable in part to its low toxicity LD50 at 2.5 g/kg (rats, oral).
See also
[edit]
Dithionite
References
[edit]
^ Weinrach, J. B.; Meyer, D. R.; Guy, J. T.; Michalski, P. E.; Carter, K. L.; Grubisha, D. S.; Bennett, D. W. (1992). "A structural study of sodium dithionite and its ephemeral dihydrate: A new conformation for the dithionite ion". Journal of Crystallographic and Spectroscopic Research. 22 (3): 291–301. doi:10.1007/BF01199531. S2CID 97124638.
^ Greenwood, Norman N.; Earnshaw, Alan (1997). Chemistry of the Elements (2nd ed.). Butterworth-Heinemann. doi:10.1016/C2009-0-30414-6. ISBN 978-0-08-037941-8.
^ a b Catherine E. Housecroft; Alan G. Sharpe (2008). "Chapter 16: The group 16 elements". Inorganic Chemistry, 3rd Edition. Pearson. p. 520. ISBN 978-0-13-175553-6.
^ a b José Jiménez Barberá; Adolf Metzger; Manfred Wolf (15 June 2000). "Sulfites, Thiosulfates, and Dithionites". Ullmann's Encyclopedia of Industrial Chemistry. Wiley Online Library. doi:10.1002/14356007.a25_477. ISBN 978-3527306732.
^ Mayhew, S. G. (2008). "The Redox Potential of Dithionite and SO−2 from Equilibrium Reactions with Flavodoxins, Methyl Viologen and Hydrogen plus Hydrogenase". European Journal of Biochemistry. 85 (2): 535–547. doi:10.1111/j.1432-1033.1978.tb12269.x. PMID 648533.
^ J. Org. Chem., 1980, 45 (21), pp 4126–4129,
^ "Aldehyde sulfoxylate systemic fungicides". google.com. Archived from the original on 27 April 2018. Retrieved 27 April 2018.
^ Božič, Mojca; Kokol, Vanja (2008). "Ecological alternatives to the reduction and oxidation processes in dyeing with vat and sulphur dyes". Dyes and Pigments. 76 (2): 299–309. doi:10.1016/j.dyepig.2006.05.041.
^ "The Best Rust Removers for Restoring Every Surface". 23 March 2023.
^ MAYHEW, Stephen G. (1978). "The Redox Potential of Dithionite and SO-2 from Equilibrium Reactions with Flavodoxins, Methyl Viologen and Hydrogen plus Hydrogenase". European Journal of Biochemistry. 85 (2): 535–547. doi:10.1111/j.1432-1033.1978.tb12269.x. ISSN 0014-2956. PMID 648533.
^ Kenneth L. Williamson "Reduction of Indigo: Sodium Hydrosulfite as a Reducing Agent" J. Chem. Educ., 1989, volume 66, p 359. doi:10.1021/ed066p359.2
^ Knight, David W.; Hartung, Jens (15 September 2006). "1-Hydroxypyridine-2(1H)-thione". 1-Hydroxypyridine-2(1H)-thione. Encyclopedia of Reagents for Organic Synthesis. John Wiley & Sons. doi:10.1002/047084289X.rh067.pub2. ISBN 978-0471936237.
^ "Kodak Direct Positive Film 5246" (PDF). 125px.com. Kodak. Retrieved 6 November 2019.
External links
[edit]
Sodium dithionite - ipcs inchem
| v t e Sodium compounds |
| Inorganic | | | | --- | | Halides | NaF NaHF2 NaCl NaBr NaI NaAt | | Chalcogenides | NaO3 NaO2 Na2O Na2O2 NaOH NaOD Na2S NaSH Na2Se NaSeH Na2Te NaHTe Na2Po | | Pnictogenides | Na3N NaN3 NaNH2 NaN2H3 Na3P Na3As Na3Sb | | Oxyhalides | NaClO NaClO2 NaClO3 NaClO4 NaBrO NaBrO2 NaBrO3 NaBrO4 NaIO NaIO3 NaIO4 | | Oxychalcogenides | Na2SO3 Na2SO4 NaHSO3 NaHSO4 Na2S2O3 Na2S2O4 Na2S2O5 Na2S2O6 Na2S2O7 Na2S2O8 Na2S4O6 Na2SeO3 Na2SeO4 NaHSeO3 Na2TeO3 | | Oxypnictogenides | NaNO2 NaNO3 Na2N2O2 Na2N2O3 NaH2PO4 NaPO3 Na2HPO4 NaPO2H2 Na2HPO3 Na2PO3F Na3PS2O2 Na3PO4 Na5P3O10 Na4P2O7 Na2H2P2O7 Na3AsO3 Na3AsO4 Na2HAsO4 NaH2AsO4 NaSbO3 | | Metalates | NaAlH4 NaAlO2 Na3AlF6 NaAl(SO4)2 NaAuCl4 Na2TiF6 NaBiO3 NaCoO2 NaMnO4 NaReO4 NaTcO3 NaTcO4 NaVO3 Na2CrO4 Na2Cr2O7 Na2Cr3O10 Na2[Fe(CO)4] Na2MnO4 Na2MoO4 Na3IrCl6 Na2PtCl6 Na2TiO3 Na2U2O7 Na2WO4 Na2Zn(OH)4 Na3VO4 Na6V10O28 Na4Fe(CN)6 Na3Fe(CN)6 Na3Fe(C2O4)3 Na3[Co(NO2)6] Na2PdCl4 | | Others | NaSbF6 NaAsF6 NaBH4 NaBH3(CN) NaBO2 Na2B4O7 Na2B2O9 Na2B8O13 NaCN NaCNO NaH NaHCO3 Na4XeO6 NaHXeO4 NaOCN NaSCN Na2CO3 Na2C2O4 Na2C3S5 Na2GeO3 Na2He Na2xSi3yO2y+x (Na2SiO3 Na6Si2O7 Na4SiO4) Na2SiF6 NaNSi2(CH3)6 | |
| Organic | CH3ONa C2H5ONa HCOONa C2H5COONa C3H7COONa Na2C4H4O6 C4H5NaO6 NaCH3COO NaC6H5CO2 NaC6H4(OH)CO2 NaC12H23O2 NaC10H8 Na2[Fe[CN5]NO] C6H16AlNaO4 NaC6H7O6 C5H8NO4Na C6H5Na C4H9Na NaC5H5 C15H31COONa C17H33COONa C18H35O2Na C164H256O68S2Na2 |
^ "Sodium dithionite - ipcs inchem" (PDF). www.inchem.org. Berliln, Germany. 2004. Archived from the original (PDF) on 17 April 2018. Retrieved 15 June 2018.
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Practical Organic Chemistry
A student handbook of techniques
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Table of contents (6 chapters)
Front Matter
Introduction
Carrying out reactions
Isolation and purification of reaction products
Separation of organic mixtures by chromatography
Preparation of samples for spectroscopy
Finding chemical information
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J. T. Sharp,
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Book Title: Practical Organic Chemistry
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Softcover ISBN: 978-0-412-28230-0Published: 26 January 1989
eBook ISBN: 978-94-009-0819-2Published: 06 December 2012
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13717 | https://brainly.com/question/27945579 | [FREE] What is \cos(28^\circ)? - brainly.com
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The trigonometric function gives the ratio of different sides of a right-angle triangle. The value of cos(28°) is -0.9626.
What are Trigonometric functions?
The trigonometric function gives the ratio of different sides of a right-angle triangle.
Sin θ=Hypotenuse PerpendicularCos θ=Hypotenuse BaseTan θ=Base PerpendicularCosec θ=Perpendicular HypotenuseSec θ=Base HypotenuseCot θ=Perpendicular Base
where perpendicular is the side of the triangle which is opposite to the angle, and the hypotenuse is the longest side of the triangle which is opposite to the 90° angle.
The cosine is the ratio of the base of the triangle and the hypotenuse of the triangle. Therefore, The value of cos(28°) is -0.9626.
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The value of cos(2 8∘) is approximately 0.8829. This value is derived from the definition of the cosine function in right triangles. You can find this using a calculator or trigonometric tables.
Explanation
To find the value of cos(2 8∘), we refer to the cosine function, which is a fundamental trigonometric function. The cosine of an angle in a right triangle is defined as the ratio of the length of the adjacent side to the length of the hypotenuse.
Understanding Cosine:
The cosine function can be understood in the context of a right-angled triangle, where:
cos(θ)=Hypotenuse Adjacent
Finding the Value:
Specifically, cos(2 8∘) can be accurately calculated using a scientific calculator set to degree mode or by using trigonometric tables.
Approximate Value:
The approximate value of cos(2 8∘) is around 0.8829.
This means that when you take the cosine of 28 degrees, you find that the ratio of the adjacent side to the hypotenuse in the corresponding right triangle measures approximately 0.8829.
Examples & Evidence
For example, in a right triangle where one angle is 28 degrees, if the length of the hypotenuse is 10 units, the length of the adjacent side can be calculated as 10×cos(2 8∘)≈8.829 units.
Using trigonometric functions established in mathematics, particularly the cosine function, confirms this value. Numerous scientific calculators and trigonometric references list cos(2 8∘) as approximately 0.8829.
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To find cos(28 degrees), you can use a calculator to get the approximate value of 0.8829. Ensure your calculator is in degree mode before performing the calculation. This value is useful in solving problems involving right-angled triangles.
The cosine function, denoted as cos(θ), is a trigonometric function representing the adjacent side divided by the hypotenuse in a right-angled triangle.
To find the cosine of 28 degrees, you can use a calculator or a trigonometric table. The value of cos(28 degrees) is approximately 0.8829.
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13718 | https://chem.libretexts.org/Bookshelves/Analytical_Chemistry/Supplemental_Modules_(Analytical_Chemistry)/Electrochemistry/Electrodes/Standard_Hydrogen_Electrode | Published Time: 2013-10-02T00:43:02Z
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Standard Electrodes
Last updated Aug 29, 2023
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Electrodes
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Table of contents
1. What are the mechanics of an electrode?
1. What processes are going on?
2. What is an electrode made of?
Examples of Electrodes
Standard Hydrogen Electrode
What is a SHE made of?
What is going on in this process?
Three Electrode System
References
Problems
Contributors and Attributions
An electrode by definition is a point where current enters and leaves the electrolyte. When the current leaves the electrodes it is known as the cathode and when the current enters it is known as the anode. Electrodes are vital components of electrochemical cells. They transport produced electrons from one half-cell to another, which produce an electrical charge. This charge is based off a standard electrode system (SHE) with a reference potential of 0 volts and serves as a medium for any cell potential calculation.
What are the mechanics of an electrode?
What processes are going on?
An electrode is a metal whose surface serves as the location where oxidation-reduction equilibrium is established between the metal and what is in the solution. The electrode can either be an anode or a cathode. An anode receives current or electrons from the electrolyte mixture, thus becoming oxidized. When the atoms or molecules get close enough to the surface of the electrode, the solution in which the electrode is placed into, donates electrons. This causes the atoms/molecules to become positive ions.
The opposite occurs with the cathode. Here the electrons are released from the electrode and the solution around it is reduced.
What is an electrode made of?
An electrode has to be a good electrical conductor so it is usually a metal. Now what this metal is made out of is dependent on whether or not it is involved in the reaction. Some reactions require an inert electrode that does not participate. An example of this would be platinum in the SHE reaction(described later). While other reactions utilize solid forms of the reactants, making them the electrodes. An example of this type of cell would be:
(left side is the anode) Cu(s)|Cu(NO 3)2(aq) (0.1M)||AgNO 3(aq) (.01M)|Ag(s) (right side is cathode)
(In the above cell set up: the outside components are the electrodes for the reaction while the inner parts are the solutions they are immersed in)
Here you can see that a solid form of the reactant, copper, is used. The copper, as well as the silver, are participating as reactants and electrodes.
Examples of Electrodes
Some commonly used inert electrodes: graphite(carbon), platinum, gold, and rhodium.
Some commonly used reactive (or involved) electrodes: copper, zinc, lead, and silver.
Standard Hydrogen Electrode
A Standard Hydrogen Electrode (SHE) is an electrode that scientists use for reference on all half-cell potential reactions. The value of the standard electrode potential is zero, which forms the basis one needs to calculate cell potentials using different electrodes or different concentrations. It is important to have this common reference electrode just as it is important for the International Bureau of Weights and Measures to keep a sealed piece of metal that is used to reference the S.I. Kilogram.
What is a SHE made of?
SHE is composed of a 1.0 M H+(aq) solution containing a square piece of platinized platinum (connected to a platinum wire where electrons can be exchanged) inside a tube. During the reaction, hydrogen gas is then passed through the tube and into the solution causing the reaction:
2H+(aq) + 2e-<==> H 2(g).
Platinum is used because it is inert and does not react much with hydrogen.
What is going on in this process?
First an initial discharge allows electrons to fill into the highest occupied energy level of Pt. As this is done, some of the H+ ions form H 3 O+ ions with the water molecules in the solution. These hydrogen and hydronium ions then get close enough to the Pt electrode (on the platinized surface of this electrode) to where a hydrogen is attracted to the electrons in the metal and forms a hydrogen atom. Then these combine with other hydrogen atoms to create H2(g). This hydrogen gas is released from the system. In order to keep the reaction going, the electrode requires a constant flow of H 2(g). The Pt wire is connected to a similar electrode in which the opposite process is occurring, thus producing a charge that is referenced at 0 volts. Other standard electrodes are usually preferred because the SHE can be a difficult electrode to set up. The difficulty arises in the preparation of the platinized surface and in controlling the concentration of the reactants. For this reason the SHE is referred to as a hypothetical electrode.
Three Electrode System
The three electrode system is made up of the working electrode, reference electrode, and the auxiliary electrode. The three electrode system is important in voltammetry. All three of these electrodes serve a unique roll in the three electrode system. A reference electrode refers to an electrode that has an established electrode potential. In an electrochemical cell, the reference electrode can be used as a half cell. When the reference electrode acts as a half cell, the other half cell's electrode potential can be discovered. An auxiliary electrode is an electrode makes sure that current does not pass through the reference cell. It makes sure the current is equal to that of the working electrode's current. The working electrode is the electrode that transports electrons to and from the substances that are present. Some examples of reference cells include:
Calomel electrode: This reference electrode consists of a mercury and mercury-chloride molecules. This electrode can be relatively easier to make and maintain compared to the SHE. It is composed of a solid paste of Hg 2 Cl 2 and liquid elemental mercury attached to a rod that is immersed in a saturated KCl solution. It is necessary to have the solution saturated because this allows for the activity to be fixed by the potassium chloride and the voltage to be lower and closer to the SHE. This saturated solution allows for the exchange of chlorine ions to take place. All this is usually placed inside a tube that has a porous salt bridge to allow the electrons to flow back through and complete the circuit.
1 2 H g 2 C l 2(s)+e−⇌H g(l)+C l−(a q)1 2 H g 2 C l 2(s)+e−⇌H g(l)+C l(a q)−
Silver-Silver Chloride electrode: An electrode of this sort precipitates a salt in the solution that participates in the electrode reaction. This electrode consists, of solid silver and its precipitated salt AgCl. This a widely used reference electrode because it is inexpensive and not as toxic as the Calomel electrode that contains mercury. A Silver-Silver Chloride electrode is made by taking a wire of solid silver and coding it in AgCl. Then it is placed in a tube of KCl and AgCl solution. This allows ions to be formed (and the opposite) as electrons flow in and out of the electrode system.
A g C l(s)+e−⇌A g+(a q)+C l−(a q)A g C l(s)+e−⇌A g(a q)++C l(a q)−
References
Ives, David J. G., and George John. Janz. "2. The Hydrogen Electrode." Reference Electrodes. New York [usw.]: Acad. Pr., 1961. Print.
Allmand, A., and Harold Johann Thomas. Ellingham. "Chapter 4: The Electrolysis Bath." The Principles of Applied Electrochemistry,. New York: Longmans, Green, 1924. Print
The Standard Hydrogen Electrode: A misrepresented concept,
Problems
Which electrode oxidizes the solution in the half-cell? Anode or Cathode?
Why is the Standard Hydrogen Electrode important to calculating cell potentials?
Identify the which side is the cathode and which side is the anode.
Ag(s) | Ag+(aq)(.5M) || Ag+(aq) (.05M) | Ag(s)
Why is it important to use an inert electrode in situations like the SHE?
What is the standard half cell potential for the SHE?
Answers (highlight to see):
Anode
It is important in calculating half-cell potentials because it serves as a reference. Without this electrode, there would be no basis to calculate values of cell potentials.
The left is the anode and the right is the cathode.
It is important to use an inert electrode in this situation because it will not react or participate in the reaction in the cell, just provide a surface area for the reaction to occur.
5.0 volts.
Contributors and Attributions
Michael Devengenzo (UCD)
Standard Electrodes is shared under a CC BY-NC-SA 4.0 license and was authored, remixed, and/or curated by LibreTexts.
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If the line y=x−8 y=x−8 intersects the parabola y 2=4 x y 2=4 x at A and B, what's the chord's length?
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How to find the coordinates of the points A and B? If one gets that, then distance formula can be used to find the length.
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edited Mar 4, 2017 at 9:21
S.C.B.
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Let us denote the intersection of y 2=4 x y 2=4 x and y=x−8 y=x−8 as (a,b)(a,b). Note that the coordinates (a,b)(a,b) satisfies
b 2=4 a=(a−8)2⟺a 2−20 a+64=(a−4)(a−16)=0 b 2=4 a=(a−8)2⟺a 2−20 a+64=(a−4)(a−16)=0
As (a−8)2−4 a=a 2−20 a+64(a−8)2−4 a=a 2−20 a+64. So the intersections are (4,−4)(4,−4) and (16,8)(16,8). So the distance between the two points is 12 2–√12 2.
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edited Mar 4, 2017 at 9:35
answered Mar 4, 2017 at 9:15
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Any point on y 2=4 x y 2=4 x can be set as (t 2,2 t)(t 2,2 t)
x−y+8=0⟹t 2−2 t−8=0⟹t 1,t 2 x−y+8=0⟹t 2−2 t−8=0⟹t 1,t 2 are 4,−2 4,−2
Now we need (t 2 1−t 2 2)2+(2 t 1−2 t 2)2−−−−−−−−−−−−−−−−−−−√(t 1 2−t 2 2)2+(2 t 1−2 t 2)2
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answered Mar 4, 2017 at 9:16
lab bhattacharjeelab bhattacharjee
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1What is the length of the intercept made by the parabola y 2=4 x y 2=4 x on the line x+y=3 x+y=3?
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Electric Field Lines (AQA A Level Physics): Revision Note
Exam code: 7408
Author
Last updated
Representing Electric Fields
Field lines are used to represent the direction and magnitude of an electric field
In an electric field, field lines are always directed from the positive charge to the negative charge
In a uniform electric field, the field lines are equally spaced at all points, this means that
The electric field strength is constant at all points in the field
The force on a test charge has the same magnitude and direction at all points in the field
In a radial electric field, the field lines spread out with distance, this means that
The field lines are equally spaced as they exit the surface of the charge
However, the separation between the field lines increases with distance
Therefore, the magnitude of electric field strength and the force on a test charge decreases with distance
Electric Field around a Point Charge
Around a point charge, the electric field is radial and the lines are directly radially inwards or outwards:
If the charge is positive(+), the field lines are radially outwards
If the charge is negative(-), the field lines are radially inwards
Electric field lines around a point charge are directed away from a positive charge and towards a negative charge
A radial field spreads uniformly to or from the charge in all directions, but the strength of the field decreases with distance
The electric field is stronger where the lines are closer together
The electric field is weaker where the lines are further apart
This shares many similarities to radial gravitational field lines around a point mass
Since gravity is only an attractive force, the field lines will look similar to the negative point charge, whilst electric field lines can be in either direction
Electric Field around a Conducting Sphere
When a conducting sphere becomes charged, the resulting electric field around it is the same as it would be if all the charge was concentrated at the centre
This means that a charged sphere can be treated in the same way as a point charge in calculations
Electric field lines around a charged conducting sphere are similar to the field lines around a point charge
Electric Field between Two Point Charges
For two opposite charges:
The field lines are directed from the positive charge to the negative charge
The closer the charges are brought together, the stronger the attractive electric force between them becomes
The electric field lines between two opposite charges are directed from the positive to the negative charge. The field lines connect the surfaces of the charges to represent attraction
For two charges of the same type:
The field lines are directed away from two positive charges or towards two negative charges
The closer the charges are brought together, the stronger the repulsive electric force between them becomes
There is a neutral point at the midpoint between the charges where the resultant electric force is zero
The electric field lines between two like charges are directed away from positive charges or towards negative charges. The field lines do not connect the surfaces of the charges to represent repulsion
Electric Field between Two Parallel Plates
When a potential difference is applied between two parallel plates, they become charged
The electric field between the plates is uniform
The electric field beyond the edges of the plates is non-uniform
Electric field lines between two parallel plates are directed from the positive to the negative plate. A uniform electric field has equally spaced field lines
Electric Field between a Point Charge and Parallel Plate
The field around a point charge travelling between two parallel plates combines
The field around a point charge
The field between two parallel plates
The electric field lines between a point charge and a parallel plate are similar to the field between two opposite charges. The field lines become parallel when they touch the plate
Worked Example
Sketch the electric field lines between the two point charges in the diagram below.
Answer:
Electric field lines around point charges have arrows which point radially outwards for positive charges and radially inwards for negative charges
Arrows (representing force on a positive test charge) point from the positive charge to the negative charge
Examiner Tips and Tricks
Always label the arrows on the field lines! The lines must also touch the surface of the source charge or plates and they must never cross.
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13721 | https://jackwestin.com/resources/mcat-content/demographic-shifts-and-social-change/urbanization | Evaluate Your Competitiveness with Free Admissions Consultation Book Now
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MCAT Content / Demographic Shifts And Social Change / Urbanization
Urbanization
Topic: Demographic Shifts And Social Change
Urbanization is the process of a population shift from rural areas to cities, often motivated by economic factors.
Urbanization is the process of a population shift from rural areas to cities. During the last century, global populations have urbanized rapidly:
13% of people lived in urban environments in the year 1900
29% of people lived in urban environments in the year 1950
Another term for urbanization is “rural flight. ” In modern times, this flight often occurs in a region following the industrialization of agriculture—when fewer people are needed to bring the same amount of agricultural output to market—and related agricultural services and industries are consolidated. These factors negatively affect the economy of small- and middle-sized farms and strongly reduce the size of the rural labor market. Rural flight is exacerbated when the population decline leads to the loss of rural services (such as business enterprises and schools), which leads to greater loss of population as people leave to seek those features.
Urbanization occurs naturally from individual and corporate efforts to reduce time and expense in commuting, while improving opportunities for jobs, education, housing, entertainment, and transportation. Living in cities permits individuals and families to take advantage of the opportunities of proximity, diversity, and marketplace competition. Due to their high populations, urban areas can also have more diverse social communities than rural areas, allowing others to find people like them.
Cities are dynamic places—they grow, shrink, and change. The fact of urban growth is undeniable: throughout the twentieth century, cities have grown rapidly. The growth machine theory of urban growth says urban growth is driven by a coalition of interest groups who all benefit from continuous growth and expansion. In some cases, that growth has been poorly controlled, resulting in a phenomenon known as urban sprawl. Urban sprawl entails the growth of a city into low-density and auto-dependent rural land, high segregation of land use (e.g., retail sections placed far from residential areas, often in large shopping malls or retail complexes), and design features that encourage car dependency.
Recently in developed countries, sociologists have observed suburbanization and counterurbanization, or movement away from cities. These patterns may be driven by transportation infrastructure, or social factors like racism. In developed countries, people are able to move out of cities while still maintaining many of the advantages of city life (for instance, improved communications and means of transportation). In fact, counterurbanization appears most common among the middle and upper classes who can afford to buy their own homes.
Some have suggested that suburbanization and urban sprawl is driven by consumer preference; people prefer to live in lower density, quieter, more private communities that they perceive as safer and more relaxed than urban neighborhoods. Such preferences echo a common strain of criticism of urban life, which tends to focus on urban decline (also referred to as urban decay). According to these critics, urban decay is caused by the excessive density and crowding of cities, and it drives out residents, creating the conditions for urban sprawl.
As cities evolve from manufacturing-based industrial to service- and information-based postindustrial societies, gentrification becomes more common. Gentrification occurs when members of the middle and upper classes enter and renovate city areas that have been historically less affluent while the poor urban underclass are forced by resulting price pressures to leave those neighborhoods for increasingly decaying portions of the city.
As cities evolved from manufacturing-based industrial to service- and information-based postindustrial societies, gentrification became the next evolution of a cities lifespan. Gentrification occurs when members of the middle and upper classes enter and renovate city areas that have been historically less affluent while the poor urban underclass is forced by resulting price pressures to leave those neighborhoods for increasingly decaying portions of the city.
Cities have responded to urban decay and urban sprawl by launching urban renewal programs. Urban renewal programs attempt to counter urban decay and restore growth as well as to make cities more pleasant and livable.
Practice Questions
Khan Academy
MCAT Official Prep (AAMC)
Practice Exam 1 P/S Section Question 30
Key Points
• Urbanization may be driven by local and global economic and social changes, and is generally a product of modernization and industrialization.
• Urbanization has economic and environmental effects. Economically, urbanization drives up prices, especially real estate, which can force original residents to move to less-desirable neighborhoods.
• Recently in developed countries, sociologists have observed suburbanization and counterurbanization, or movement away from cities, which may be driven by transportation infrastructure, or social factors like racism.
• Suburbanization may be driven by white flight.
• Counterurbanization refers, broadly, to movement away from the city, which may include urban-to-rural migration and suburbanization.
• Counterurbanization has created shrinking cities and attempts to better control urban growth.
• Urban sprawl results when cities grow uncontrolled, expanding into rural land and making walking, public transit, or bicycling impractical.
• Critics of urban life often focus on urban decay, which may be self-perpetuating, according to the broken windows theory.
• Urban renewal attempts to counter urban decay and restore growth.
Key Terms
Suburbanization: A term used to describe the growth of areas on the fringes of major cities; one of the many causes of the increase in urban sprawl.
Rural flight: A term used to describe the migratory patterns of peoples from rural areas into urban areas.
Industrialization: The development of industries in a country or region on a wide scale.
Urbanization: The physical growth of urban areas as a result of rural migration and even suburban concentration into cities.
Counterurbanization: A demographic and social process whereby people move from urban areas to rural areas.
Gentrification: A shift in an urban community toward wealthier residents and/or businesses and increasing property values; often resulting in poorer residents being displaced by wealthier newcomers.
White flight: The large-scale migration of whites of various European ancestries, from racially mixed urban regions to more racially homogeneous suburban areas.
Counterurbanization: Counterurbanisation is a demographic and social process whereby people move from urban areas to rural areas.
Urban decline/Urban Decay: The deterioration of the inner city often caused by lack of investment and maintenance.
Urban growth:An increase in the urbanized land cover.
Urban renewal: Urban renewal refers to programs of land redevelopment in areas of moderate- to high-density urban land use.
Urban sprawl: Refers to the unrestricted growth in many urban areas of housing, commercial development, and roads over large expanses of land, with little concern for urban planning.
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13722 | https://brainly.in/question/15916877 | Write down the decimal expansions of the following rational numbers by writing their denominators in the - Brainly.in
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Priscillaselvam703
11.03.2020
Math
Secondary School
answered • expert verified
Write down the decimal expansions of the following rational numbers by writing their denominators in the form 2ᵐ × 5ⁿ, where m, n are non-negative integers.(i) 3/8 (ii) 13/125
(iii) 7/80 (iv) 14588/625
(v) 129/(2² × 5⁷)
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nikitasingh79
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Concept :
If the factors of denominator of the given rational number is of form 2ᵐ × 5ⁿ,where n and m are non negative integers, then the decimal expansion of the rational number is terminating otherwise non terminating recurring.
SOLUTION :
(i) Given : 3/8
We can write it as : 3/ 2³
Here, the factors of the denominator 8 are 2³× 5^0, which is in the form 2ᵐ × 5ⁿ .
So , 23/8 has terminating decimal expansion.
The decimal expansion of 3/ 8 :
⅜ = 3×5³ / 2³ × 5³
= 3 × 125 / (2×5)³
= 375 / 10³ = 375 /1000
= 0.375
Hence, the decimal expansion of 3/8 is 0.375.
(ii) Given : 13/125
We can write it as :13/5³
Here, the factors of the denominator 125 are 2⁰ × 5³ which is in the form 2ᵐ × 5ⁿ .
So , 13/125 has a terminating decimal expansion.
The decimal expansion of 13/ 125 :
13/125 = 13 ×2³ /5³× 2³
= 13 × 8 /(5×2)³
= 104 / 10³ = 104/1000
= 0.104
Hence, the decimal expansion of 13/125 is 0.104.
(iii) Given : 7/80
Here, the factors of the denominator 80 are 2⁴ × 5¹, which is in the form 2ᵐ × 5ⁿ .
So , 7/80 has a terminating decimal expansion.
The decimal expansion of 7/80 :
7/80 = 7 × 5³ / 2⁴ × 5¹ × 5³
= 7 × 125 / (2×5)⁴
= 875/10⁴ = 875/10000
= 0.0875
Hence, the decimal expansion of 7/80 is 0.0875 .
(iv) Given : 14588/625
Here, the factors of the denominator 625 is 2⁰ × 5⁴, which is in the form 2ᵐ 5ⁿ
So , 14588/625 has a terminating decimal expansion.
The decimal expansion of 14588/625 :
14588/625 = 14588 × 2⁴ / 5⁴× 2⁴
= 14588 × 16 / (5×2)⁴
= 233408/ 10⁴
= 23.3408
Hence, the decimal expansion of 14588/625 is 23.3408.
(v) Given : 129/(2² × 5⁷)
129/(2² × 5⁷)
Here, the factors of the denominator are 2² × 5⁷ which is in the form 2ᵐ × 5ⁿ .
So ,129/2² × 5⁷ has a terminating decimal expansion.
The decimal expansion of 129/2² × 5⁷ :
129 × 2⁵ /2² × 5⁷ × 2⁵
= 129 × 32 / (2 × 5)⁷
= 4128 / 10000000
= 0.0004128
Hence, the decimal expansion of 129/(2² × 5⁷) is 0.0004128.
HOPE THIS ANSWER WILL HELP YOU…
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Without actually performing the long division, state whether state whether the following rational numbers will have a terminating decimal expansion or a non-terminating repeating decimal expansion.
(i)
(ii)
(iii)
(iv)
(v)
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13723 | https://www.youtube.com/watch?v=39JulN5g0S4 | 05 Percentiles and the Normal Distribution
ProfThacker
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Transcript:
okay we're going to finish this section with a talk on percentiles and the applications of percentiles we did talk about percentiles earlier in the semester um in chapter 3 we looked at percentiles and we discovered that the percentile rank of a data value X indicates the percent of data values within a distribution that are less than the data value X so let's look at a real example um if for example we're talking about data value X representing um children's Heights um let's suppose for let's suppose we're looking at a certain age group and we're told that the percentile rank for 48 in is 90% what that means is that 90% of all the other Heights in the that distribution are less than 48 in in a normal distribution the percentile rank of data value X is equal to the percent of area under the normal curve to the left or below the data value X so suppose that we have raw scores labeled on the horizontal axis from least to greatest and suppose data value X or maybe 48 in whatever raw data we're representing is labeled right here the percent of area to the left of that raw score represents the percentile rank of that raw score it's basically the percent of values within the distribution that are less than that raw score and that is the definition of a percentile rank we're going to always express a percentile rank as a whole number so this is one of the different things about a percentile ranks when we find a proportion of area on our calculator we know that it gives us a proportion as a as a decimal less than one to turn it into a whole number we're going to first multiply it by 100 to turn it into a percent and then we're going to round it to the nearest whole number so let's look at an example Matthew earned a grade of 87 on his History exam if the grades in his class were normally distributed so this is a this is what makes this question different than a chapter 3 question if we're told it's normally distributed then we we know that we can use normal CDF in this problem the mean is 80 and the standard deviation is s find Matthew use percentile rank on the exam so to find the percentile rank we simply need to find the area to the left of Matthew's score the area to the left of Matthew's score will give us the proportion of of class grades that were less than Matthews so the percentile rank of data value x equal 87 is equal to the percent of area under the normal curve to the left which means less than or below the data value x equal 87 so in our normal distribution we're told that the mean is 80 the standard deviation is seven mean is 80 standard deviation is seven Matthew scored more than the mean so his value is on the right side of the mean so we make a vertical line representing the point where Matthew's score is and the area to the left of Matthew's score is basically the percent of students or percent of grades that were less than 87 if we can take that percent of area turn it into a uh and round it to the nearest whole number that will give us the percentile rank for 87 the percentile rank of area will be the percent of area to the left of 87 and round to the nearest whole number so we go to our calculator to find the percent of area to the left of 87 so we should note that the lower score that we enter in our calculator is this is incorrect sorry just should be negative 100,000 for um negative Infinity the upper score is positive 87 the mean is 80 the standard deviation is 7 and we get an area of 8413 so let me just correct this real quick this should be negative 100,000 okay so negative um sorry the proportion of area is 8413 to turn this into a percentile rank we simply have to turn it into a percent 84.1 3% and then round to the nearest whole number which is 84 so we're going to follow normal rounding rules round to the nearest whole number so to put into words what we just found is is the percentile rank for Matthew's grade 87 is 84 it's a little hard when you when you put it in this when you phrase it like this it's hard sometimes to recognize which one is the score and which one is the percentile rank but you'll get used to it with practice 87 is the score and approximately 84% of all the grades fall fall at or below 87 so it's percentile rank is is 84 we could also use percentile notation P with a subscript 84 means the percentile rank um is 84 for the data value 87 so data value 87 has a percentile rank that is 84 okay what this means is if the percentile rank of 87 is 84 then the score 87 is the 84th percentile the score 87 which is Matthew's score is the 84th percentile which what we know from chapter 3 means that about 84% of grades were at or below 87 approximately 84% of the students had scores less than or equal to the score 887 try example 7.10 Miss Brook took two exams last week her grades and class results were as follows okay on English she got an 83 the class meme was 80 the standard deviation was six on math she got a 7 7even the class mean was 74 the standard deviation was five assuming the results of both exams to be normally distributed which basically tells us that this is different from chapter 3 because we're using the normal distribution and therefore the bell curve and normal CDF used percentile ranks to determine on which exam Miss Brook did better relative to her class so we're going to calculate the percent rank for each of her grades and once we find the percentile rank the one that is higher will be the one in which Miss Brooke did better relative to her class okay so if you want pause the video find the two percentile ranks for 83 and then for for 77 and then decide which one she did better relative to her class to find the percentile rank of the English exam we calculate the area to the left of the raw score 83 so this is what our picture looks like in the center we have a mean of 80 and a standard deviation of six 83 is to the right of 80 so we label it here we make our vertical line and we shade all the area to the left of 83 because that is the that is the meaning of a percentile rank it's the area to the left of the score so on our calculator to find the area to the left of this score we have to use normal CDF with a negative uh with a with a lower value of negative 100,000 oh again I messed this up with a lower value of 100,000 and upper value of 83 a mean of 80 and a standard deviation of six okay so putting those four values in you should get a value of 6915 which is 69.1 15% which converts into a percentile rank of 69 the percentile rank of 83 is 69 it means that approximately 69% of the students scored at or below an 83 now to find find the percentile rank of the math exam we calculate the area to the left of the raw score 77 so this is a different picture we have a different mean and standard deviation because we're looking at a different distribution of scores this is the math class so mean of 74 a standard deviation of five 77 is on the right side of 74 so we make a vertical line and we shade all the area to the right of um to the left of 74 sorry to the left of 77 this area to the left of 77 represents the percent of area that is or the percent of grades that are less than 77 so we're going to go to our calculator to find out what that actual percent is the lower value is 00,000 the upper value is 77 the mean is 74 the standard deviation is five and so we get 7257 on our calculator which converts to 72.5 s% that is about a 73 when rounded to the nearest number so the score 77 has a percentile rank of 73 which means that about 73% of the students scored at or below a 77 so now let's recap our findings the percentile rank of 69 for the English grade means that 69% of the score of of the class scored at or below Miss Brook's score the percentile rank of 73 for the math exam means that 73% of the class scored at or below Miss Brook's score so looking at which one has the higher percentile rank 69 or 73 73 for the math grade means that the math grade because it has the higher percentile rank is the grade that's better relative to the rest of the class okay so that concludes the lecture on um percentiles and percentile ranks so what you want to take away from this lecture is that percentile rank is found by looking at the area to the left of a particular raw score so you're going to start at negative Infinity then go up to the raw score and then put in the mean and the standard deviation and that will give you a proportion of area you'll convert that proportion of area to a percent and then round to the nearest whole number and that will give you the percentile rank any percentile rank that is higher is the one that is better relative to the rest of the class okay so we've done we've discovered how to use infv Norm today and we did a little bit more with normal CDF using real raw scores and real world applications so if you have any questions please email me I'll see you at the next lecture |
13724 | https://blog.csdn.net/qinghuaci666/article/details/82056055 | 博客
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数列求和之差比数列
DashingPuma的博客
12-13
705
Q:an=n⋅3n−1,求Sn.Q:a_{n}=n\cdot 3^{n-1},求S_{n}.Q:an=n⋅3n−1,求Sn.
A:Sn=(n2−14)⋅3n+14.A:S_{n}=\left( \dfrac {n}{2}-\dfrac {1}{4}\right) \cdot 3^{n}+\dfrac {1}{4}.A:Sn=(2n−41)⋅3n+41.
...
等差数列与等比数列
weixin_33949359的博客
07-17
522
数列
按照一定顺序排列着的一列数称为数列
数列中的每一个数叫做这个数列的项
数列中的每一项都和它的序号有关,排在第一位的数称为这个数列的第一项(通常也叫做首项)
数列的一般形式可以写成
a1,a2,a3,…,an,…,
简记为{an}, 项数有限的数列叫做有穷数列,项数无限的数列叫做无穷数列
按照数列的每一项随序号变化的情况对数列分类:
...
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等差数列/等比数列通项公式与求和公式_等比数列求通项公式
9-25
其中等差数列的首项为a1,末项为an,项数为n,公差为d,前n项和为Sn。 二、等比数列 如:1 3 9 27 81... 通项公式: 求和公式: 其中a1为首项,q为等比数列公比,Sn为等比数列前n项和。
递归、回溯算法常用数学基础公式_递归 数学公示
9-17
等比数列 求和公式:Sn=(首项-末项公比)÷(1-公比) 通项公式: an=a1q^(n-1) 等差数列 通项公式:an = a1 + (n-1)d 求和公式:(a1+an)n/2
等差等比数列
wulong710的专栏
05-14
694
//等差数列 //an = a1+ (n-1)p; //sn = n(a1+an)/2; //等比数列 //an = a1pow(q,n-1) //sn = a1(1-pow(q,n) )/(1-q)
在线等比数列求和计算器
Linux,Java,SpringBoot,Python,Lua略知一点
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在线等比数列求和计算器
在线等比数列求和计算器
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等差等比数列公式
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本文介绍了等差数列和等比数列的基本概念及其关键计算公式,包括求和公式、项数公式及末项公式等,为理解和解决数学问题提供帮助。 摘要生成于C知道,由 DeepSeek-R1 满血版支持,前往体验 > 等差: 和=(首项+末项)×项数÷2 项数=(末项-首项)÷公差+1 ...
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本篇文章主要探讨了如何使用NumPy生成等差数列和等比数列,以及Python中实现等差数列的不同方法。 首先,NumPy提供了两种主要的方法来生成等差数列:numpy.linspace和numpy.arange。numpy.linspace函数允许...
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等比数列的应用3:如果{an}是各项均为正数的等比数列,那么{lg an}是等差数列,公差为lg q(n ≥ 2)。 习题 例1:已知{an}是等比数列且an > 0,a2a4 + 2a3a5 + a4a6 = 25,求a3 + a5的值。 例2:有四个数,...
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也就是对于首项为1,公比为2的等比数列中,前n-1项之和(2^n-1)加一等于第n项(2^n)。根据满二叉树可得,an就是叶子结点数,s(n-1)就是有两个儿子的节点数。 我们再以此推广到完全二叉树:每一个正常结点变成一个单儿子节点,就会损失相同个数的叶子节点和有两个儿子的节点(包括他自己;如果是在上图的情况...
等差、等比数列公式总结.pdf
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【标题】"vb课设 等差等比数列求和"涉及到的是使用Visual Basic (VB)编程语言实现一个课程设计项目,该项目的核心功能是计算等差数列和等比数列的和。在VB中,我们可以创建图形用户界面(GUI)应用程序,通过用户...
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2.fx在负无穷到正无穷上的积分等于0 等比数列是指从第二项起,每一项与它的前一项的比值等于同一个常数的一种数列,常用G、P表示。这个常数叫做等比数列的公比,公比通常用字母q表示(q≠0),等比数列a1≠ 0。其中{an}中的每一项均不为0。注:q=1 时,an为常数列。
第九周-OJ-D等比数列
qq_36260328的博客
11-01
441
问题及代码:
/
Copyright(c)2016,烟台大学计算机学院
All rights reserved.
文件名称:tset.cpp
作者:张旺
完成日期:2016年10月22日
版本号:v1.0
题目描述
已知q与n,求等比数列之和: 1+q+q2+q3+q4+…+qn
输入
输入数据含有不多于50对的数据,每对数据含有一个整数n(1≤n≤20),一个小数q(0<q<
差比数列通解(1.6)
qq_38445477的博客
08-07
803
错位相减法
如果差位为n次?
如何用程序解决(迭代)
include
include
include
int jiecheng(int n)
{
if(n<=1){return 1;}
if(n>1){return njiecheng(n-1);}
}
int zuhe(int n,int m)
{
if(m>=n){return 1;}
return jiecheng(n)
等差,等比和差比数列推导
RAVANLALA
07-04
1937
先讲讲等差数列
比如我们有一个数列1,2,3,4,5,6,7,8.......,n,很明显公差d = 1,首项a1 = 1.
所以an = a1 + (n - 1)d ;
又可以推导出来an = am + (n - m)d; (已知数列的第m项 和d求出第n项)
我们的等差数列的前n项和为a1 + a2 +a3 + ...... + an = sn
从an = a1 + (n - 1)d ...
等差数列和等比数列的公式、法则、定理
kwame211的博客
12-01
1万+
一、 等差数列
如果一个数列从第二项起,每一项与它的前一项的差等于同一个常数,这个数列就叫做等差数列,这个常数叫做等差数列的公差,公差常用字母d表示.
等差数列的通项公式为:
an=a1+(n-1)d (1)
前n项和公式为:
Sn=na1+n(n-1)d/2或Sn=n(a1+an)/2(2)
从(1)式可以看出,an是n的一次数函(d≠0)或常数函数(d=0),(n,an)排在一
等差、等比数列求和公式
hyqsong的专栏
08-03
9943
等比数列求和公式求和公式推导(1)Sn=a1+a2+a3+…+an(公比为q)
(2)qSn=a1q+a2q+a3q+…+anq=a2+a3+a4+…+an+a(n+1)
(3)Sn-qSn=(1-q)Sn=a1-a(n+1)
(4)a(n+1)=a1q^n
(5)Sn=a1(1-q^n)/(1-q)(q≠1)等差数列求和公式Sn=(A1+An)n/2
图解等差数列和等比数列求和公式
pathfinder1987的博客
06-17
1万+
1. 等差数列求和公式2. 等比数列求和公式
二叉树中序遍历和后序遍历的递归与非递归算法
Tseng
03-09
6797
昨天写的前序遍历的递归与非递归算法,在非递归算法中主要还是借用到了栈这一工具,其实在中序遍历和后序遍历中依旧可以理由栈的特性来进行非递归的遍历 操作。
1.中序遍历
1.1 中序遍历的递归算法
二叉树的构造以在上一篇博文中有说到,这里就不再重复了,直接上代码
public void InOrderTraverse(TreeNode T){
if(T==null){
等比数列前N项和的公式推导
cv小学dl三年级
04-10
3368
设等比数列的前n项和为S(n), 等比数列的第一项为a1,比值为q。
(1)S(n) = a1 + a1 q + a1 q ^ 2 + .... + a1 q ^ (n - 1);(2)S(n+1) = a1 + a1 q + a1 q ^ 2 + .... + a1 q ^ (n - 1) + a1 q ^ n;由(2)式减(1)式得(3)S(n+1) - s(n)
伯努利分布(二项分布)的假设检验
热门推荐
呼呼的博客
08-14
3万+
要点
1. 单个二项分布检验用SPSS二项检验或者单样本T检验
2. 比较两个个二项分布差异性之类的可以用Anova或者独立样本T检验,后者可以得到置信区间。
3.上面的SPSS数据形式都是1和0的形式,用MATLAB生成即可。
譬如这样一个问题:
中国的互联网络覆盖率是不是在30%以上(5%显著性水平)?抽样显示,150个样本中,有57个是有网络覆盖...
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13725 | https://www.numdam.org/item/10.5802/aif.3094.pdf | A N N A L E S D E L ’ I N S T I T U T F O U R I E R ANNALES DE L’INSTITUT FOURIER Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER On evil Kronecker sequences and lacunary trigonometric products Tome 67, no 2 (2017), p. 637-687.
© Association des Annales de l’institut Fourier, 2017, Certains droits réservés.
Cet article est mis à disposition selon les termes de la licence CREATIVE COMMONS ATTRIBUTION – PAS DE MODIFICATION 3.0 FRANCE.
L’accès aux articles de la revue « Annales de l’institut Fourier » ( implique l’accord avec les conditions générales d’utilisation (
cedram Article mis en ligne dans le cadre du Centre de diffusion des revues académiques de mathématiques Ann. Inst. Fourier, Grenoble 67, 2 (2017) 637-687 ON EVIL KRONECKER SEQUENCES AND LACUNARY TRIGONOMETRIC PRODUCTS by Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER () Abstract. — An important result of Weyl states that for every sequence (nk)k⩾1 of distinct positive integers the sequence of fractional parts of (nkα)k⩾1 is u.d. mod 1 for almost all α. However, in this general case it is usually ex-tremely difficult to measure the speed of convergence of the empirical distribution of ({n1α}, . . . , {nNα}) towards the uniform distribution. In this paper we investi-gate the case when (nk)k⩾1 is the sequence of evil numbers, that is the sequence of non-negative integers having an even sum of digits in base 2. We utilize a connec-tion with lacunary trigonometric products QL ℓ=0 sin π2ℓα , and by giving sharp metric estimates for such products we derive sharp metric estimates for exponen-tial sums of (nkα)k⩾1 and for the discrepancy of ({nkα})k⩾1 . Furthermore, we provide some explicit examples of numbers α for which we can give estimates for the discrepancy of ({nkα})k⩾1.
Résumé. — Un résultat important de Weyl nous dit que pour chaque suite (nk)k⩾1 de nombres entiers positifs différents la suite {nkα}k⩾1 est équidistribuée modulo 1 pour presque tous les réels α. Dans ce cas, il est d’habitude extrêmement difficile de mesurer la vitesse de convergence de la distribution empirique vers l’équidistribution.
Dans cet article, nous étudions le cas ou (nk)k⩾1 est la suite des nombres entiers « méchants », donc la suite des nombres positifs la une somme de chiffres paire dans la base 2. Nous relions ce probléme aux produits trigonométriques QL l=0 ∥sin π2lα∥ en donnant des estimations exactes pour de tels produits et nous obtenons des estimations exactes pour la discrépance de la suite {nkα}k⩾1.
En plus, nous donnons des exemples concrets de réels α pour lesquels nous pouvons obtenir des estimations pour la discrépance de la suite {nkα}k⩾1.
Keywords: evil numbers, Thue–Morse sequence, (nα)-sequence, discrepancy, lacunary trigonometric products.
Math. classification: 11B85, 11K38, 11B83, 11A63, 68R15.
() The first and the third author are supported by the Project I 1751-N26 “Multiplica-tivity, Determinism, and Randomness”, as well as by the Austrian Science Fund (FWF): Project F5507-N26, which is part of the Special Research Program “Quasi-Monte Carlo Methods: Theory and Applications”. The second author is supported by the Austrian Science Fund (FWF): Project F5505-N26, which is part of the Special Research Program “Quasi-Monte Carlo Methods: Theory and Applications”.
638 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER 1. Introduction and statement of results Throughout the rest of this paper, we let s2(n) denote the sum-of-digits function of n in base 2, and we let (nk)k⩾1 denote the sequence of non-negative integers which have an even sum-of-digits function in base 2, sorted in increasing order. Furthermore, we write (mk)k⩾1 for the sequences of those numbers which are not contained in (nk)k⩾1, sorted in increasing or-der; thus (mk)k⩾1 = (1, 2, 4, 7, 8, 11, . . . ). The numbers (nk)k⩾1 are called evil numbers, while the numbers (mk)k⩾1 are called odious numbers. These sequences are connected to the well-known Thue–Morse sequence (tn)n⩾0 = (0, 1, 1, 0, 1, 0, 0, 1, 1, 0, 0, 1, 0, 1, 1, 0, . . .), a binary sequence which is also de-fined in terms of s2(n); more precisely, tn = 1 if s2(n) is odd and tn = 0 otherwise. The Thue–Morse sequence has many remarkable properties; for an extensive survey see . In our notation we have n ∈(nk)k⩾1 if and only if tn = 0.
In this paper we analyze exponential sums of the form PN k=1 e2πinkα for reals α ∈[0, 1), and – what is intimately connected – products of the form QL ℓ=0 sin π2ℓα , as well as distribution properties of the sequence ({nkα})k⩾1. To quantify the regularity of the distribution of a finite set of real numbers in [0, 1] we use the notion of the star-discrepancy D∗ N. For given numbers x1, . . . , xN, their star-discrepancy is defined by D∗ N(x1, . . . , xN) = sup a∈[0,1] 1 N N X n=1 10,a −a .
An infinite sequence (xn)n⩾1 whose discrepancy D∗ N tends to zero as N → ∞is called uniformly distributed modulo one (u.d. mod 1). Informally speaking, the star-discrepancy is a measure for the deviation between the uniform distribution on [0, 1] and the empirical distribution of a given point set; in probabilistic terminology this corresponds to the Kolmogorov– Smirnov statistic. Discrepancy theory is a rich subject, which has close links to number theory, probability theory, ergodic theory and numerical analy-sis. For more information on discrepancy theory, we refer to the standard monographs [12, 27].
Sequences of the form ({nα})n⩾1 are called Kronecker sequences. One of the fundamental results of discrepancy theory states that such a sequence is u.d. mod 1 if and only if α is irrational. It is also well-known that the discrepancy of such a sequence depends on Diophantine approximation ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 639 properties of α. More precisely, we have Ω m(N) X n=1 an = ND∗ N({α}, . . . , {Nα}) = O m(N) X n=1 an as N →∞, where a1, a2, a3, . . . are the continued fraction coefficients of α and where m(N) is defined by qm(N)−1 < N ⩽qm(N) with q1 < q2 < q3 < . . . denoting the best approximation denominators of α (see for example [12, Corollary 1.64]). Hence ND∗ N({α}, . . . , {Nα}) = O (log N) as N →∞ if α has bounded continued fraction coefficients, and, as a consequence of metric results of Khintchine , for every ε > 0 we have ND∗ N({α}, . . . , {Nα}) = O (log N)(log log N)1+ε) N as N →∞ for almost all α ∈R.
It is known (see ) and it will be re-proved implicitly in this paper (see Section 6) that the sequence ({nkα})k⩾1, which we will call the evil Kronecker sequence, is also uniformly distributed in the unit interval if and only if α is irrational. However, it turns out to be a very difficult task to give sharp estimates for the discrepancy of this sequence for concrete values of α.
As we will see, the discrepancy of an evil Kronecker sequence ({nkα})k⩾1 depends on Diophantine approximation properties and properties of the digit representation of α in base 2. Until now there are only few (non-trivial) cases of α where we have enough information about both of these aspects.
Exponential sums and discrepancy theory are intimately connected. One such connection is Weyl’s criterion, two others are the Erdős–Turán in-equality and Koksma’s inequality. The Erdős–Turán inequality (see for ex-ample [12, 27, 38]) states that for points x1, . . . , xN ∈[0, 1] we have D∗ N(x1, . . . , xN) ⩽ 1 H + 1 + H X h=1 1 h 1 N N X k=1 e2πihxk , (1.1) where H is an arbitrary positive integer. Koksma’s inequality says that (1.2) Z 1 0 f(x) dx −1 N N X k=1 f(xk) ⩽(Var[0,1]f)D∗ N(x1, . . . , xN), TOME 67 (2017), FASCICULE 2 640 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER for any function f having bounded variation Var[0,1]f on [0, 1]. When com-bined, the Erdős–Turán inequality and Koksma’s inequality show that ex-ponential sums can be used to obtain both upper and lower bounds for the discrepancy.
As an explicit lower bound from (1.2) we get (compare for example ): (1.3) D∗ N (x1, . . . , xN) ⩾ 1 4H 1 N N X k=1 e2πiHxk , where H is an arbitrary positive integer. Koksma’s inequality and its multi-dimensional generalization are also the cornerstone of the application of low-discrepancy point sets in numerical integration (so-called Quasi-Monte Carlo integration; see for example [11, 34]).
Consequently, in this paper we will mainly be concerned with the problem of investigating exponential sums of the form PN k=1 e2πinkα. It turns out that this investigation relies on studying lacunary products of the form QL ℓ=0 sin π2ℓα . Furthermore we study the discrepancy of ({nkα})k⩾1. For all three topics we obtain sharp metric results. The investigation of the lower bound for the discrepancy leads to a challenging open problem in Diophantine approximation. Finally, we consider two concrete non-trivial special examples for α.
The main results of this paper are the following. (Throughout the rest of this paper, we write exp(x) for ex.) Theorem 1.1. — Let (nk)k⩾1 be the sequence of evil numbers, sorted in increasing order, and let h ̸= 0 be an integer. Let ε > 0 be arbitrary.
Then for almost all α ∈(0, 1) we have (1.4) N X k=1 e2πihnkα ⩽exp π √log 2 + ε (log N)1/2(log log log N)1/2 for all N ⩾N0(α, h, ε), and (1.5) N X k=1 e2πihnkα ⩾exp π √log 2 −ε (log N)1/2(log log log N)1/2 for infinitely many N.
Note that the exponential function in (1.4) grows more slowly than any (fixed) power of N; but faster than any (fixed) power of log N. In other words, as a consequence of Theorem 1.1 for every ε > 0 and every A > 0 ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 641 we have Ω (log N)A = N X k=1 e2πihnkα = O (N ε) as N →∞, for almost all α.
It will turn out that Theorem 1.1 is an almost immediate consequence of the following result on lacunary trigonometric products.
Theorem 1.2. — Let ε > 0 be arbitrary. Then for almost all α ∈(0, 1) we have (1.6) L Y ℓ=0 2 sin π2ℓα ⩽exp (π + ε) p L log log L for all L ⩾L0 (α, ε), and (1.7) L Y ℓ=0 2 sin π2ℓα ⩾exp (π −ε) p L log log L for infinitely many L.
This result is a consequence of a more general result, Theorem 2.5, which will be formulated later in Section 2 since it needs some technical prereq-uisites.
From the lower bound in Theorem 1.1 and formula (1.3) we immediately obtain a metric lower bound for the discrepancy D∗ N ({n1α} , . . . , {nNα}) of this sequence. However in Theorem 1.3 it turns out that the true metric order of the discrepancy D∗ N ({n1α} , . . . , {nNα}) is much larger.
Theorem 1.3. — Let (nk)k⩾1 be the sequence of evil numbers, sorted in increasing order. Let ε > 0 be arbitrary. Then for almost all α ∈(0, 1) we have ND∗ N ({n1α} , . . . , {nNα}) = O N 1+ log λ log 2 +ε as N →∞, (1.8) and ND∗ N ({n1α} , . . . , {nNα}) ⩾ N 1+ log λ log 2 −ε as N →∞ (1.9) for infinitely many N. Here λ is a real constant defined below for which it is known that (1.10) 0.66130 < λ < 0.66135.
TOME 67 (2017), FASCICULE 2 642 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER The number λ in Theorem 1.3 appears in a result of Fouvry and Mauduit , which states that (1.11) I1(L) := Z 1 0 L−1 Y ℓ=0 sin π2ℓα dα = κλL (1 + o(1)) for L →∞, with constants κ > 0 and λ with 0.654336 < λ < 0.663197. In Lemma 5.1 in Section 5 we will improve the estimate for λ to (1.10). Note that as a consequence of (1.8) and (1.10) we have ND∗ N ({n1α} , . . . , {nNα}) = O N 0.404 as N →∞, for almost all α. This should be compared with the general metric discrep-ancy bound (1.12) ND∗ N ({b1α} , . . . , {bNα}) = O √ N(log N)3/2+ε as N →∞ for almost all α, which holds for every strictly increasing sequence of pos-itive integers (bk)k⩾1 (see ). It is known that in the general setting the upper bound given by (1.12) is optimal (up to powers of logarithms; see ).
Thus the upper bound given in Theorem 1.3 is significantly stronger than the general metric discrepancy bound given by (1.12). Furthermore we want to emphasize the fact that the precision of Theorem 1.3 is quite remarkable, in view of the fact that good bounds for the typical order of the discrepancy are only known for a very small number of classes of parametric sequences.
One of the main objectives of Theorems 1.1 and 1.3 is to examine the degree of “pseudorandomness” of the parametric sequences ({nkα})k⩾1, and consequently also of the sequence of evil numbers and the Thue–Morse sequence. By classical probability theory, for a sequence X1, X2, . . . of in-dependent, identically distributed (i.i.d.) random variables having uniform distribution on [0, 1] we have the law of the iterated logarithm (LIL) (1.13) lim sup N→∞ PN k=1 e2πihXk √2N log log N = 1 √ 2 almost surely (a.s.) and the Chung–Smirnov LIL for the Kolmogorov–Smirnov statistic (that is, for the discrepancy) (1.14) lim sup N→∞ ND∗ N(X1, . . . , XN) √2N log log N = 1 2 a.s.; in other words, for a random sequence of points exponential sums are typi-cally of asymptotic order roughly √ N, and the discrepancy is typically also of the corresponding asymptotic order. Furthermore, similar results usually hold for exponential sums of (rkα)k⩾1 and for the discrepancy of ({rkα})k⩾1 when (rk)k⩾1 is a “random” increasing sequence of integers. In the simplest ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 643 model, when for every number n ⩾1 we decide independently and with fair probability whether it should be contained in (rk)k⩾1 or not, then (1.13) holds almost surely (with respect to the probability space over which the rk’s are defined) for almost all α. In a similar fashion both results (1.13) and (1.14) essentially remain valid when the random sequence (rk)k⩾1 is constructed in a more complicated fashion (see for example [19, 20, 36]).
Thus Theorem 1.1 and Theorem 1.3 show that almost surely the asymp-totic order of exponential sums and of the discrepancy of ({nkα})k⩾1 for the evil numbers (nk)k⩾1 does not match the corresponding order in the random case, by this means showing an interesting deviation from “pseu-dorandom” behavior of the sequence (nk)k⩾1 itself. On the other hand, the behavior of ({nkα})k⩾1 also does not match the behavior of nα-sequences for typical values of α. More precisely, as already mentioned above, as a consequence of metric results of Khintchine and due to the fact that the discrepancy of ({nα})n⩾1 can be expressed in terms of the continued fraction expansion of α, we have (1.15) D∗ N({α}, {2α}, . . . , {Nα}) = O (log N)(log log N)1+ε) N as N →∞for almost all α. Consequently, by Theorem 1.3, the typical asymptotic order of the discrepancy of parametric sequences ({nkα})k⩾1 is significantly larger than that of typical nα-sequences, and by Theorem 1.1 this is also true for exponential sums. Thus, with respect to exponen-tial sums as well as with respect to the discrepancy, parametric sequences (nkα)k⩾1 generated by the evil numbers (nk)k⩾1 occupy a position some-where between nα-sequences and truly random sequences. We also want to comment on the fact that there is a huge difference between the order of the exponential sums in Theorem 1.1 and the order of the discrepancy in Theorem 1.3. This is a very surprising phenomenon, which is related to problems from metric Diophantine approximation (which are implicit in the proof of Theorem 1.3, and are briefly discussed in the concluding Section 7).
As already mentioned earlier, it is rather difficult to give the right order for the exponential sums in Theorem 1.1, the trigonometric products in Theorem 1.2, and the discrepancy of ({nkα})k⩾1 for concrete non-trivial examples of α. What do we mean by a “non-trivial” example? In the first part of Section 6 we will point out the following facts: • The order of the discrepancy of the pure Kronecker sequence ({nα})n⩾1 is never significantly larger than the order of the dis-crepancy of the evil Kronecker sequences ({nkα})k⩾1.
TOME 67 (2017), FASCICULE 2 644 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER • If the order of the discrepancy D∗ N of the pure Kronecker sequence satisfies ND∗ N = Ω N log 3 log 4 then the discrepancy e D∗ N of the evil Kronecker sequences is essentially of the same order as the discrep-ancy of the pure Kronecker sequence.
• If the order of the discrepancy D∗ N of the pure Kronecker sequence satisfies ND∗ N = O N log 3 log 4 then e D∗ N satisfies N e D∗ N = O N log 3 log 4 +ε .
Thus an “interesting non-trivial” example means for us an example where α is a “natural” real number such as √ 2, e, π (it seems to us that there is no chance to handle these numbers since we do not have enough information on their digit representation), or where α is such that e D∗ N and hence D∗ N is small (say N e D∗ N = O (N ε) – however we cannot give such examples) or where the quality of the distribution of the sequences ({nα})n⩾1 and ({nkα})k⩾1 differ strongly. Two such examples are given in Theorem 1.4.
Especially in the first example the difference between D∗ N and e D∗ N is of the maximal possible form.
Theorem 1.4.
(a) Let α = 2 3 +P∞ k=1 1 42k . Then for the star-discrepancy D∗ N of the pure Kronecker sequence ({nα})n⩾1 we have ND∗ N = O (log N) , whereas for the star-discrepancy e D∗ N of the evil Kronecker sequence ({nkα})k⩾1 we have N e D∗ N = O N log 3 log 4 +ε and N e D∗ N = Ω N log 3 log 4 −ε for every ε > 0.
(b) Let γ = 0.1001011001101001 . . . be the Thue–Morse real in base 2.
Then for the star-discrepancy D∗ N of the pure Kronecker sequence ({nγ})n⩾1 we have ND∗ N = O (N ε) for all ε > 0, whereas for the star-discrepancy e D∗ N of the evil Kronecker sequence ({nkγ})k⩾1 we have N e D∗ N = Ω N 0.6178775 .
We would like to point out here that there is an intimate connection be-tween distribution properties of ({nkα})k⩾1 and of certain types of hybrid ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 645 sequences. For some information on the analysis of hybrid sequences see for example , and .
As already mentioned, the proofs of Theorems 1.1 and 1.3 are based on a connection between exponential sums of (nkα)k⩾1 and the lacunary trigonometric products studied in Theorem 1.2. We will establish this con-nection in the following lines, and exploit it in Section 2 in more detail.
For the time being, we assume that N is of the form 2L for some positive integer L.
To analyze the exponential sums appearing in Theorem 1.1 and on the right-hand side of (1.1), we define Sh(N) = N X k=1 e2πihnkα.
By the assumption that N = 2L we have Sh(N) = 2N−1 X n=0 1 2 X δ∈{0,1} exp 2πiδ s2(n) 2 exp (2πihnα) = 1 2 X δ∈{0,1} X (η0,...,ηL)∈{0,1}L+1 exp 2πi L X ℓ=0 δ ηℓ 2 + ηℓh2ℓα !
= 1 2 X δ∈{0,1} L Y ℓ=0 X ηℓ∈{0,1} exp 2πiηℓ δ 2 + h2ℓα , which yields |Sh(N)| ⩽1 2 L Y ℓ=0 |2 sin πh2ℓα| + 1 2 L Y ℓ=0 |2 cos πh2ℓα| (1.16) and |Sh(N)| ⩾ 1 2 L Y ℓ=0 |2 sin πh2ℓα| −1 2 L Y ℓ=0 |2 cos πh2ℓα| .
(1.17) A similar analysis for the sequence (mk)k⩾1 shows that (1.18) N X k=1 e2πihmkα ⩽1 2 L Y ℓ=0 |2 sin πh2ℓα| + 1 2 L Y ℓ=0 |2 cos πh2ℓα|, where again we assume that N = 2L.
By taking logarithms, we can convert the trigonometric products appear-ing in (1.16) and (1.18) into so-called lacunary sums; these sums have been TOME 67 (2017), FASCICULE 2 646 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER intensively investigated in Fourier analysis, and a wide range of mathemat-ical methods is available for studying them (see the following Section 2).
Thus the theory of lacunary sums allows us to obtain an estimate for the size of the exponential sums Sh(N) in the case when N is a power of 2; however, it will turn out that we may also drop the condition that N is an integral power of 2 by applying a dyadic decomposition method.
Note that by (1.15) and by the fact that the evil numbers have asymptotic density 1/2 it is easy to show that all the conclusions of Theorem 1.1 and Theorem 1.3 remain valid if we replace the sequence (nk)k⩾1 by the sequence of odious numbers (mk)k⩾1.
The outline of the remaining part of this paper is as follows. In Section 2 we explain the main principles of the theory of lacunary (trigonometric) sums, and state several lemmas as well as Theorem 2.5, which we require for the proofs of Theorem 1.1 and Theorem 1.2. In Section 3 we give the proofs for the results stated in Section 2, and in Section 4 we give the proofs of Theorem 1.1 and Theorem 1.2. In Section 5, we prove Theorem 1.3, and in Section 6 we prove Theorem 1.4. Finally, in Section 7, we briefly mention a problem from metric Diophantine approximation, which was posed by LeVeque in and is related to the proof of Theorem 1.3.
2. Probabilistic results for lacunary trigonometric products It is a well-known fact that so-called lacunary systems of trigonometric functions, that is, systems of the form (cos 2πsℓα)ℓ⩾1 or (sin 2πsℓα)ℓ⩾1 for rapidly increasing (sℓ)ℓ⩾1, exhibit properties which are typical of sequences of independent random variables. This similarity includes the central limit theorem, the law of the iterated logarithm, and Kolmogorov’s “Three series” convergence theorem. The situation is particularly well understood when (sℓ)ℓ⩾1 satisfies the Hadamard gap condition (2.1) sℓ+1 sℓ ⩾q > 1, ℓ⩾1.
To a certain degree this almost-independence property extends to systems (f(sℓα))ℓ⩾1 for a function f which is periodic with period one and satisfies certain regularity properties; however, in this case the number-theoretic properties of (sℓ)ℓ⩾1 play an important role, and the almost-independent behavior generally fails when (2.1) is relaxed to a weaker growth condi-tion. The case which has been investigated in the greatest detail is that when f has bounded variation on [0, 1], since this case is (by Koksma’s ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 647 inequality) closely connected to the discrepancy of the sequence of frac-tional parts ({sℓα})ℓ⩾1, which in turn can be interpreted as the (one-sided) Kolmogorov–Smirnov statistic adapted to the case of the uniform measure on [0, 1].
To estimate the trigonometric products appearing in (1.16) and (1.18) we will use the equalities L−1 Y ℓ=0 2 sin πh2ℓα = exp L−1 X ℓ=0 log 2 sin πh2ℓα !
(2.2) and L−1 Y ℓ=0 2 cos πh2ℓα = exp L−1 X ℓ=0 log 2 cos πh2ℓα !
, (2.3) respectively, to transform the problem of lacunary trigonometric products into a problem concerning lacunary sums. However, the functions f1(α) := log |2 sin πα| (2.4) and f2(α) := log |2 cos πα| (2.5) do not have bounded variation in [0, 1] (see the figures below). Conse-quently, the known results are not applicable in this situation, and we have to adopt the proof techniques in such a way that they can handle this kind of problem.(1) 0.2 0.4 0.6 0.8 1.0 -2.5 -2.0 -1.5 -1.0 -0.5 0.5 0.2 0.4 0.6 0.8 1.0 -2.5 -2.0 -1.5 -1.0 -0.5 0.5 Figure 2.1. The functions f1 (left) and f2 (right).
(1) There exist a few results concerning lacunary series when f is neither required to have bounded variation, nor to be Lipschitz- or Hölder-continuous, nor to have a modulus of continuity of a certain regularity; see for example . However, for these results the growth requirements for (sℓ)ℓ⩾1 are much stronger than (2.1), which means that they are not applicable in our case, since by (2.2) and (2.3) we have to deal with lacunary sequences growing exactly with the speed presumed in (2.1), and not faster.
TOME 67 (2017), FASCICULE 2 648 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER In the following we will assume that f is a measurable real function satisfying (2.6) f(α + 1) = f(α), Z 1 0 f(α) dα = 0, Z 1 0 f(α)2 dα < ∞, and that (sℓ)ℓ⩾1 is a sequence satisfying (2.1). We write f(α) ∼ ∞ X j=1 aj cos 2πjα + bj sin 2πjα for the Fourier series of f, and we will assume that the Fourier coefficients of f satisfy (2.7) |aj| ⩽1 j , |bj| ⩽1 j , for j ⩾1.
The inequalities in line (2.7) appear frequently in the theory of lacunary series, since the upper bound stated there describes precisely (up to mul-tiplication with a constant) the maximal asymptotic order of the Fourier coefficients of a function of bounded variation (see for example [41, p. 48]).
However, even if the function is not of bounded variation the estimates in (2.7) may still be true; this can be seen by the fact that for the (un-bounded) functions f1 and f2 from lines (2.4) and (2.5), respectively, we have f1(α) ∼ ∞ X j=1 −1 j cos 2πjα and f2(α) ∼ ∞ X j=1 (−1)j+1 j cos 2πjα.
By the way, we note that both f1 and f2 are even functions, and that both of them satisfy (2.6).
Throughout the remaining part of this paper, we will write P for the Lebesgue measure on the unit interval. Note that the unit interval, equipped with Borel sets and Lebesgue measure, is a probability space, and that accordingly every measurable function on [0, 1] can be seen as a random variable over this probability space. We also write ∥· ∥2 for the L2(0, 1) norm and ∥· ∥∞for the supremum norm of a function, respectively.
The main technical tool in this section is the following exponential in-equality (Lemma 2.1). Together with the subsequent lemmas it will allow us to give an upper bound for the measure of those α for which (2.2) and (2.3) are large (stated in Lemma 2.4). We state Lemmas 2.1–2.4 in a slightly more general form than necessary for the proofs of Theorem 1.1 ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 649 and 1.2, and we will use them to prove an additional new theorem, namely Theorem 2.5 below.
Lemma 2.1. — Assume that f is an even measurable function satisfy-ing (2.6), whose Fourier coefficients satisfy (2.7). Furthermore, let (sℓ)ℓ⩾1 be a sequence of positive integers satisfying (2.1) for some number q > 1.
Then there exists a number L0 = L0(q), such that the following holds. Let L ⩾L0 be given, and write p(α) for the L8-th partial sum of the Fourier series of f. Then for all (2.8) λ ∈ h 0, L−1/9i we have Z 1 0 exp λ L X ℓ=1 p(sℓα) !
dα ⩽exp 2λ2π2L 3 and Z 1 0 exp λ L X ℓ=1 p(sℓα) !
dα ⩽exp 12λ2 √q √q −1∥p∥1/2 2 L .
The same two conclusions hold if f is an odd function instead of an even function.
We emphasize the fact that the number L0 in the statement of Lemma 2.1 depends only on the growth parameter q; it does not depend on the function f or the sequence (sℓ)ℓ⩾1. The same will be true for the numbers L0(q) in Lemmas 2.2 and 2.4 below.
From Lemma 2.1 we will deduce the following Lemma 2.2, which is a large-deviation bound for the maximal partial sum of a lacunary sums.
Lemma 2.2. — Let f and (sℓ)ℓ⩾1 be as in Lemma 2.1. Then there exists a number L0 = L0(q), such that the following holds. Let L ⩾L0 be given, and assume that L is an integral power of 2. Write p(α) for the L8-th partial sum of the Fourier series of f. Then we have P α ∈(0, 1) : max 1⩽M⩽L M X ℓ=1 p(sℓα) > 29q q −1 p L log log L !
⩽ 48 (log L)1.4 and, under the additional assumption that ∥p∥1/4 2 ⩾L−1/100, we also have P α ∈(0, 1) : max 1⩽M⩽L M X ℓ=1 p(sℓα) > 43∥p∥1/4 2 √q √q −1 p L log log L + √ L !
⩽ 45 (log L)2 .
TOME 67 (2017), FASCICULE 2 650 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER The same two conclusions hold if f is an odd function instead of an even function.
Lemma 2.3. — Let f and (sℓ)ℓ⩾1 be as in Lemma 2.1. Let L be given, and write r(α) for the remainder term of the L8-th partial sum of the Fourier series of f. Then we have Z 1 0 max 1⩽M⩽L M X ℓ=1 r(sℓα) !2 dα ⩽4 L4 From Lemmas 2.2 and 2.3 we will deduce the following Lemma 2.4.
Lemma 2.4. — Let f and (sℓ)ℓ⩾1 be as in Lemma 2.1. Then there exists a number L0 = L0(q), such that the following holds. Let L ⩾L0 be given, and assume that L is an integral power of 2. Then we have P α ∈(0, 1) : max 1⩽M⩽L M X ℓ=1 f(sℓα) > 30q q −1 p L log log L !
⩽ 49 (log L)1.4 and, under the additional assumption that ∥p∥1/4 2 ⩾L−1/100, we also have P α ∈(0, 1) : max 1⩽M⩽L M X ℓ=1 p(sℓα) > 43∥p∥1/4 2 √q √q −1 p L log log L + 2 √ L !
⩽ 46 (log L)2 .
The same two conclusions hold if f is an odd function instead of an even function.
As a consequence of Lemma 2.4 we obtain the following theorem, which is a bounded law of the iterated logarithm and is of some interest in its own right. As far as we know, this is the first law of the iterated logarithm for Hadamard lacunary function series which can be applied to a class of unbounded functions f.
Theorem 2.5. — Assume that f is an even measurable function satis-fying (2.6), whose Fourier coefficients satisfy (2.7). Furthermore, let (sℓ)ℓ⩾1 be a sequence of positive integers satisfying (2.1) for some number q > 1.
Then we have lim sup L→∞ PL ℓ=1 f(sℓα) √L log log L ⩽cq for almost all α ∈(0, 1), where we can choose cq = max 85q q −1, 122∥f∥1/4 2 √q √q −1 .
ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 651 We note in passing that from our proofs it seems that the conclusion of Theorem 2.5 remains true if the conditions |aj| ⩽j−1, |bj| ⩽j−1 in (2.7) are relaxed to |aj| ⩽j−1/2−ε, |bj| ⩽j−1/2−ε for some fixed ε > 0; however, in this case the constant cq has to be replaced by some other constant cq,ε which may also depend on ε. We will not pursue this possible generalization any further in the present paper.
For the proofs of Theorem 1.1 and 1.2 we will also need the following result. It has first been stated by Fortet ; a concise proof can be found in . This result can be seen as a special case of the more general results in .
Lemma 2.6. — Let f be a function satisfying (2.6), which additionally satisfies a Hölder continuity condition of order β for some β > 0. Then lim sup L→∞ PL−1 ℓ=0 f(2ℓα) √2L log log L = σf for almost all α, where (2.9) σ2 f = lim m→∞ 1 m Z 1 0 f(α) + · · · + f(2m−1α) 2 dα.
3. Proofs of results from Section 2 Proof of Lemma 2.1. — The proof of Lemma 2.1, as well as the proofs of Lemmas 2.2, 2.3, 2.4 and Theorem 2.5, uses methods of Takahashi and Philipp .
Assuming that f is even, the L8-th partial sum of the Fourier series of f is of the form p(α) = L8 X j=1 aj cos 2πjα, where by assumption the coefficients aj satisfy the inequality on the left-hand side of (2.7). We note that (2.7) implies that (3.1) ∥p∥∞⩽ L8 X j=1 1 j ⩽1 + 8 log L.
We divide the set of integers {1, . . . , L} into blocks ∆1, . . . , ∆w of con-secutive numbers, for some appropriate w, such that every block contains logq 4L8 numbers (the last block may contain less).(2) More precisely, (2) We assume throughout that L is sufficiently large such that all appearing logarithms are well-defined and positive.
TOME 67 (2017), FASCICULE 2 652 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER we set w = L ⌈logq (4L8)⌉ and ∆i = n (i −1) logq 4L8 + 1, . . . , i logq(4L8) o ∩{1, . . . , L}, for 1 ⩽i ⩽w. We set I1 = Z 1 0 exp 2λ X 1⩽i⩽w, i even X ℓ∈∆i p(sℓα) dα and I2 = Z 1 0 exp 2λ X 1⩽i⩽w, i odd X ℓ∈∆i p(sℓα) dα.
Then by the Cauchy–Schwarz inequality we have (3.2) Z 1 0 exp λ L X ℓ=1 p(sℓα) !
dα ⩽ p I1I2.
Writing Ui = X ℓ∈∆i p(sℓα), 1 ⩽i ⩽w, and using the inequality ex ⩽1 + x + x2, which is valid for |x| ⩽1, we have (3.3) I1 = Z 1 0 Y 1⩽i⩽w, i even exp 2λ X ℓ∈∆i p(sℓα) dα ⩽ Z 1 0 Y 1⩽i⩽w, i even 1 + 2λUi + 4λ2U 2 i dα, where we used the fact that by (2.7), (2.8) and (3.1) we have |2λUi| ⩽2L−1/9∥p∥∞|∆i| ⩽2L−1/9(1 + 8 log L) logq(4L8) ⩽1 for L ⩾L0(q) ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 653 Using the classical trigonometric identity (3.4) (cos y)(cos z) = cos(y + z) + cos(y −z) 2 , for y, z ∈R, we have U 2 i = X ℓ∈∆i L8 X j=1 aj cos 2πjsℓα 2 = X ℓ1,ℓ2∈∆i X 1⩽j1,j2⩽L8 aj1aj2 cos(2π(j1sℓ1 + j2sℓ2)) 2 (3.5) +cos(2π(j1sℓ1 −j2sℓ2)) 2 =: Vi + Wi.
(3.6) Here we write Vi for the sum of all those cosine-functions having frequencies in the interval smin(∆i), 2L8smax(∆i) , where min(∆i) and max(∆i) denote the smallest resp. largest element of ∆i, and we write Wi for the sum of those cosine-functions having frequencies smaller than smin(∆i). It is easy to check that no other frequencies can occur in (3.5). We note that all the frequencies of the cosine-functions in Ui are also contained in the interval smin(∆i, 2L8smax(∆i) , and write (3.7) Xi = 2λUi + 4λ2Vi.
Using this notation we have Y 1⩽i⩽w, i even 1 + 2λUi + 4λ2U 2 i = Y 1⩽i⩽w, i even 1 + Xi + 4λ2Wi .
From Minkowski’s inequality and (2.7) we deduce that Wi ⩽1 2 X ℓ1,ℓ2∈∆i X 1⩽j1,j2⩽L8 | {z } |j1sℓ1−j2sℓ2|j2sℓ2/sℓ1−1 |aj1aj2| (3.8) ⩽ X ℓ1,ℓ2∈∆i, ℓ1⩽ℓ2 X 1⩽j1,j2⩽L8, j1>j2sℓ2/sℓ1−1 1 j1j2 TOME 67 (2017), FASCICULE 2 654 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER ⩽ X ℓ1,ℓ2∈∆i, ℓ1⩽ℓ2 L8 X j=1 1 j 1 ⌈jsℓ2/sℓ1 −1⌉ | {z } ⩽2sℓ1/(jsℓ2) ⩽2 X ℓ1,ℓ2∈∆i, ℓ1⩽ℓ2 L8 X j=1 sℓ1 j2sℓ2 ⩽2|∆i| q q −1 π2 6 .
(3.9) Another way of continuing from line (3.8) is to use the Cauchy–Schwarz inequality, which leads to Wi ⩽ X ℓ1,ℓ2∈∆i, ℓ1⩽ℓ2 X 1⩽j2⩽L8 1 j2 X 1⩽j1⩽L8, j1>j2sℓ2/sℓ1−1 |aj1| ⩽ X ℓ1,ℓ2∈∆i, ℓ1⩽ℓ2 X 1⩽j2⩽L8 1 j2 X 1⩽j1⩽L8, j1>j2sℓ2/sℓ1−1 ⩽1/j2 1 z}|{ a2 j1 | {z } ⩽4sℓ1/(j2sℓ2) 1/2 X 1⩽j1⩽L8, j1>j2sℓ2/sℓ1−1 a2 j1 1/2 | {z } ⩽∥p∥1/2 2 ⩽2 X ℓ1,ℓ2∈∆i, ℓ1⩽ℓ2 X 1⩽j⩽L8 r sℓ1 j3sℓ2 ∥p∥1/2 2 ⩽6 √q √q −1∥p∥1/2 2 .
(3.10) Now assume that i1 < i2 are two indices from the set {1, . . . , w}, and that both i1 and i2 are even. Then by construction the frequency of any trigonometric function in Xi2 is at least twice as large as the frequency of any trigonometric function in Xi1. To see why this is the case, we recall that the frequency of the largest trigonometric function in Xi2 is at most 2L8smax(∆i1), that the frequency of the smallest trigonometric function in Xi1 is at least smin(∆i2), and that by (2.1) smin(∆i2) smax(∆i1) ⩾qmin(∆i2)−max(∆i1) = q⌈logq(4L8)⌉ ⩾4L8.
ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 655 As a consequence for every set of distinct indices i1, . . . , iv (where the cardi-nality v is arbitrary), all of which are even and are contained in {1, . . . , w}, the functions Xi1, . . . , Xiv are orthogonal, i.e., (3.11) Z 1 0 Xi1 · · · · · Xiv dα = 0 (this argument is explained in more detail in [35, 40]). Thus by (3.3), (3.6), (3.7), (3.9), and (3.11) we have I1 ⩽ Z 1 0 Y 1⩽i⩽w, i even 1 + Xi + 4λ2Wi dα ⩽ Z 1 0 Y 1⩽i⩽w, i even 1 + Xi + 4λ2|∆i|π2 3 q q −1 dα = Z 1 0 Y 1⩽i⩽w, i even 1 + 4λ2|∆i|π2 3 q q −1 | {z } does not depend on α dα ⩽ Y 1⩽i⩽w, i even exp 4λ2|∆i|π2 3 q q −1 = exp X 1⩽i⩽w, i even 4λ2|∆i|π2 3 q q −1 .
In the same way we can get an upper bound for I2, and thus by (3.2) we finally obtain Z 1 0 exp λ L X ℓ=1 p(sℓα) !
dα ⩽exp X 1⩽i⩽w, i even 2λ2|∆i|π2 3 q q −1 exp X 1⩽i⩽w, i odd 2λ2|∆i|π2 3 q q −1 = exp 2λ2Lπ2 3 q q −1 .
TOME 67 (2017), FASCICULE 2 656 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER This proves the first conclusion of the lemma. In the same way we can use (3.10) (and the corresponding upper bound for I2) to obtain Z 1 0 exp λ L X ℓ=1 p(sℓα) !
dα ⩽exp 12λ2 √q √q −1∥p∥1/2 2 , which proves the second conclusion of the lemma.
Thus we have proved both parts of Lemma 2.1 in the case when f is even; the proof in the odd case can be carried out in exactly the same way.
□ Proof of Lemma 2.2. — By assumption L is an integral power of 2. We set ν = log2 L. By classical dyadic decomposition, we can write every subset {1, . . . , M} of {1, . . . , L} as the disjoint sum of at most one set of cardinality 2ν−1, at most one set of cardinality 2ν−2, at most one set of cardinality 2ν−3, and so on, at most one set of cardinality 2⌈ν/4⌉, and additionally at most one set of cardinality at most 2⌈ν/4⌉, where all these sets contain consecutive positive integers. To be able to represent every sets {1, . . . , M} in this way, we need 2µ sets of cardinality 2ν−µ, for µ ∈{1, . . . , ν −⌈ν/4⌉}, and all the sets of cardinality at most 2⌈ν/4⌉starting at an integer multiple of 2⌈ν/4⌉. More precisely, the sets of cardinality 2ν−µ are of the form {j2ν−µ + 1, . . . , (j + 1)2ν−µ}, j ∈{0, . . . , 2µ −1}, µ ∈{1, . . . , ν −⌈ν/4⌉}, and the sets of cardinality at most 2ν/4 are of the form n j2⌈ν/4⌉+ 1, , . . . , j2⌈ν/4⌉+ w o , j ∈ n 0, . . . , 2ν−⌈ν/4⌉−1 o , w ∈ n 1, . . . , 2⌈ν/4⌉o .
For j ∈{0, . . . , 2µ −1} and µ ∈{1, . . . , ν −⌈ν/4⌉}, we set Gµ,j = α ∈(0, 1) : (j+1)2ν−µ X ℓ=j2ν−µ+1 p(sℓα) > 8q q −1 √ 2ν−µp log log 2ν−µ + µ √ 2ν−µ Using the first part of Lemma 2.1 with λ = p log log(2ν−µ) √ 2ν−µ ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 657 (note that by our construction this value of λ is admissible in Lemma 2.1, provided that L is sufficiently large) we have Z 1 0 exp p log log(2ν−µ) √ 2ν−µ (j+1)2ν−µ X ℓ=j2ν−µ+1 p(sℓα) dα ⩽exp 2 log log(2ν−µ)qπ2 3(q −1) , and consequently P(Gµ,j) ⩽exp 2π2 3 −8 q q −1 log log(2ν−µ) −µ ⩽ 1 eµ(log(2ν−µ))1.4 .
Thus we have P ν−⌈ν/4⌉ [ µ=1 2µ−1 [ j=0 Gµ,j ⩽ ν−⌈ν/4⌉ X µ=1 2µ eµ log 2⌈ν/4⌉1.4 ⩽ 20 (log L)1.4 .
(3.12) Now we set Hj = α ∈(0, 1) : max 1⩽w⩽2⌈ν/4⌉ j2⌈ν/4⌉+w X ℓ=j2⌈ν/4⌉+1 p(sℓα) > √ L , j ∈{0, . . . , 2ν−⌈ν/4⌉−1}.
By (2.7), Minkowski’s inequality, and the Carleson–Hunt inequality (see for example ), we have Z 1 0 max 1⩽w⩽2⌈ν/4⌉ j2⌈ν/4⌉+w X ℓ=j2⌈ν/4⌉+1 p(sℓα) 6 dα 1/6 ⩽ L8 X j=1 |aj| Z 1 0 max 1⩽w⩽2⌈ν/4⌉ j2⌈ν/4⌉+w X ℓ=j2⌈ν/4⌉+1 cos 2πjsℓα 6 dα 1/6 ⩽(1 + 8 log L)cabs Z 1 0 (j+1)2⌈ν/4⌉ X ℓ=j2⌈ν/4⌉+1 cos 2πjsℓα 6 dα 1/6 (3.13) for some absolute constant cabs. Estimating the integral in (3.13) can be reduced (via (3.4)) to the problem of counting the number of solutions TOME 67 (2017), FASCICULE 2 658 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER (ℓ1, . . . , ℓ6) of the Diophantine equation sℓ1 ± · · · ± sℓ6 = 0, for indices ℓ1, . . . , ℓ6 in the respective index range; we have Z 1 0 (j+1)2⌈ν/4⌉ X ℓ=j2⌈ν/4⌉+1 cos 2πjsℓα 6 dα 1/6 ⩽cq 2ν/41/2 for some constant cq depending only on q (see, for example, ). As a consequence by Markov’s inequality we obtain P(Hj) ⩽ˆ cq(1 + 8 log L)6 L6/8 L3 = ˆ cq(1 + 8 log L)6L−9/4 and P 2ν−⌈ν/4⌉−1 [ j=0 Hj ⩽L3/4ˆ cq(1 + 8 log L)6L−9/4 = ˆ cq(1 + 8 log L)6L−6/4 (3.14) for some constant ˆ cq depending only on q. We set F = ν−⌈ν/4⌉ [ µ=1 2µ−1 [ j=0 Gµ,j ∪ 2ν−⌈ν/4⌉−1 [ j=0 Hj .
Then by (3.12) and (3.14) we have P (F) ⩽ 23 (log L)1.4 + ˆ cq(1 + 8 log L)6L−6/4 (3.15) ⩽ 24 (log L)1.4 (3.16) for sufficiently large L.
By the dyadic decomposition described at the beginning of this proof, for every α ∈F C (where F C denotes the complement of the set F) we have M X ℓ=1 p(sℓα) ⩽ ν−⌈ν/4⌉ X µ=1 8q q −1 √ 2ν−µp log log 2ν−µ + µ √ 2ν−µ + √ L ⩽ 8q q −1(1 + √ 2) p L log log L + (4 + 3 √ 2) √ L + √ L ⩽29q q −1 p L log log L ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 659 for all possible values M ∈{1, . . . , L}. Thus we have proved that P α ∈(0, 1) : max 1⩽M⩽L M X ℓ=1 p(sℓα) !
29q q −1 p L log log L !
⩽ 24 (log L)1.4 (3.17) for sufficiently large L. Note that whenever the function f satisfies the assumptions of Lemma 2.1, then the function −f also satisfies these as-sumptions. Thus applying exactly the same arguments as above to the functions −f and −p instead of f and p we also obtain P α ∈(0, 1) : max 1⩽M⩽L M X ℓ=1 (−p(sℓα)) !
29q q−1 p L log log L !
⩽ 24 (log L)1.4 for sufficiently large L, which, together with (3.17), proves the first conclu-sion of Lemma 2.2.
The proof of the second conclusion of Lemma 2.2 is very similar to that of the first conclusion of the lemma. We use the same dyadic decomposition, but now we define Gµ,j = ( α ∈(0, 1) : (j+1)2ν−µ X ℓ=j2ν−µ+1 p(sℓα) > 14∥p∥1/4 2 √q √q −1 + ∥p∥1/4 2 µ × √ 2ν−µp log log 2ν−µ ) and use the second part of Lemma 2.1 with λ = ∥p∥−1/4 2 p log log(2ν−µ) √ 2ν−µ Note that this choice of λ is admissible, due to the restrictions that µ ⩽ ν −⌈ν/4⌉and ∥p∥1/4 2 ⩾L−1/100. We obtain that Z 1 0 exp ∥p∥−1/4 2 p log log(2ν−µ) √ 2ν−µ (j+1)2ν−µ X ℓ=j2ν−µ+1 p(sℓα) dα ⩽exp 12√q log log(2ν−µ) √q −1 , and consequently P(Gµ,j) ⩽exp (12 −14) √q √q −1 log log(2ν−µ) −µ ⩽ 1 eµ(log(2ν−µ))2 .
TOME 67 (2017), FASCICULE 2 660 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER The sets Hj can be defined in the same way as in the proof of the first part of the lemma. Using similar calculations we obtain that max 1⩽L⩽M M X ℓ=1 p(sℓα) ⩽ ∞ X µ=1 2−µ/2 (14 + µ) ∥p∥1/4 2 √q √q −1 p L log log L + √ L ⩽43∥p∥1/4 2 √q √q −1 p L log log L + √ L, except for a set of measure at most ∞ X µ=1 2µ eµ(log(L1/4))2 !
+ ˆ cq(1 + 8 log L)6L−9/4 ⩽ 45 (log L)2 (provided that L is sufficiently large). This proves the second part of Lemma 2.2.
□ Proof of Lemma 2.3. — The lemma follows from a simple application of Minkowski’s inequality. We have Z 1 0 max 1⩽M⩽L M X ℓ=1 r(sℓα) !
2 dα ⩽ L X M=1 Z 1 0 M X ℓ=1 r(sℓα) !
2 dα ⩽ L X M=1 M X ℓ=1 ∥r∥2 !2 .
(3.18) By (2.7) we have ∥r∥2 ⩽ ∞ X j=L8+1 2 j2 1/2 ⩽ 2 L7/2 , which implies, together with (3.18), that Z 1 0 max 1⩽M⩽L M X ℓ=1 r(sℓα) !
2 dα ⩽4 L4 .
This proves the lemma.
□ Proof of Lemma 2.4. — Assume that we have decomposed f = p + r as in Lemmas 2.1, 2.2 and 2.3. By Lemma 2.3 and Markov’s inequality we have P α ∈(0, 1) : max 1⩽M⩽L M X k=1 r(sℓα) > √ L !
⩽4 L5 .
ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 661 Together with Lemma 2.2 this yields P α ∈(0, 1) : max 1⩽M⩽L M X k=1 f(sℓα) > 29q q −1 + 1 | {z } ⩽30q/(q−1) p L log log L !
⩽ 48 (log N)1.4 + 4 L5 ⩽ 49 (log L)1.4 for sufficiently large L, which is the first part of Lemma 2.4.
In the same way we can deduce the second conclusion of Lemma 2.4 from a combination of the second conclusion of Lemma 2.2 and Lemma 2.3.
□ Proof of Theorem 2.5. — From Lemma 2.2 and Lemma 2.3 we can easily deduce Theorem 2.5, using standard methods. Let us first assume that f is either even or odd. For m ⩾1, we let Em denote the set defined by Em = ( α ∈(0, 1) : max 1⩽M⩽2m M X ℓ=1 f(sℓα) > 30q q −1 p 2m log log 2m ) .
Then by Lemma 2.6 we have (3.19) P(Em) ⩽ 49 (log 2m)1.4 which implies that ∞ X m=1 P(Em) < ∞.
Thus by the Borel–Cantelli lemma with probability 1 only finitely many events Em happen; in other words, for almost all α ∈(0, 1) we have max 1⩽M⩽2m M X ℓ=1 f(sℓα) ⩽30q q −1 p 2m log log 2m for m ⩾m0(α).
As a consequence for almost all α ∈(0, 1) we have (3.20) lim sup L→∞ PL ℓ=1 f(sℓα) √L log log L ⩽ √ 2 30q q −1, which proves Theorem 2.5 in the case when f is either even or odd. For general f we apply (3.20) to the even and odd part separately, which results in an extra multiplicative factor of 2. Note that 2 √ 2 · 30 ⩽85.
TOME 67 (2017), FASCICULE 2 662 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER In the same way we can use the second conclusion of Lemma 2.4 to obtain (3.21) lim sup L→∞ PL ℓ=1 f(sℓα) √L log log L ⩽2 √ 2 · 43 | {z } ⩽122 ∥p∥1/4 2 √q √q −1, for almost all α ∈(0, 1). Theorem 2.5 now follows from a combination of (3.20) and (3.21).
□ 4. Exponential Sums and Trigonometric Products: Proofs of Theorem 1.1 and Theorem 1.2 Theorem 1.1 and Theorem 1.2 will follow easily from the following lemma, which is a version of Lemma 2.6 in the case when the function f only sat-isfies (2.7) (instead of being Hölder-continuous). We state it only for the special case of the two functions f1 and f2 from (2.4) and (2.5).
Lemma 4.1. — For almost all α we have lim sup L→∞ PL−1 ℓ=0 f1(2ℓα) √2L log log L = π √ 2 and lim sup L→∞ PL−1 ℓ=0 f2(2ℓα) √2L log log L = 0.
Proof. — As already noted, the functions f1 and f2 satisfy condi-tions (2.6) and (2.7). Let a number d be given, and write p1 for the d-th partial sum of the Fourier series of f1, and r1 for the remainder term. Then by Lemma 2.6 we have (4.1) lim sup L→∞ PL−1 ℓ=0 p1(2ℓα) √2L log log L = σp1 for almost all α where σp1 is defined according to (2.9) (for the function p1). Note that Lemma 2.6 is applicable since p is a trigonometric polynomial (and con-sequently is also Lipschitz-continuous). Note furthermore that by (2.7) we have ∥r1∥2 ⩽ d+1 X j=1 1 j2 1/2 ⩽d−1/2.
Consequently by Theorem 2.5 we have (4.2) lim sup L→∞ PL−1 ℓ=0 r1(2ℓα) √2L log log L ⩽122d−1/8 √ 2 √ 2 −1 for almost all α.
ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 663 Replacing f1 by f2 and replacing p1 and r1 by p2 and r2, respectively, we obtain (4.1) and (4.2) with p1 and r1 replaced by p2 and r2, respectively.
Some standard calculations show that σp1 →σf1 and σp2 →σf2 as d →∞.
Furthermore, the expression on the right-hand side of (4.2) clearly tends to zero as d →∞. Thus, overall we have (4.3) lim sup L→∞ PL−1 ℓ=0 f1(2ℓα) √2L log log L = σf1 for almost all α and (4.4) lim sup L→∞ PL−1 ℓ=0 f2(2ℓα) √2L log log L = σf2 for almost all α, where σf1 and σf2 are defined according to (2.9). Calculating the values of σf1 and σf2 is a simple exercise, using the Fourier series expansion of f1 and f2, respectively; it turns out that in our specific setting we have (4.5) σ2 f1 = ∞ X j=1 1 j2 + 2 ∞ X r=1 1 2rj2 !
= π2 2 and (4.6) σ2 f2 = ∞ X j=1 1 j2 + 2 ∞ X r=1 (−1)j(−1)2rj | {z } =(−1)j 1 2rj2 = 0.
By applying the same arguments to −f1 and −f2 instead of f1 and f2 we can get absolute values in (4.3) and (4.4), if we wish. This proves Lemma 4.1.
□ Proof of Theorem 1.1 and of Theorem 1.2.
Part 1: upper bounds. — By periodicity it is obviously sufficient to prove Theorem 1.1 for h = 1. Let ε > 0 and α be given, and set ˆ ε = ε/2.
We will assume that α is an element of the set of full measure for which the conclusion of Lemma 4.1 holds. Then we have (4.7) L X ℓ=0 f1(2ℓα) ⩽ π √ 2 + ˆ ε p 2L log log L and (4.8) L X ℓ=0 f2(2ℓα) ⩽ˆ ε p 2L log log L TOME 67 (2017), FASCICULE 2 664 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER for all L ⩾L0(α). Equation (4.7) already gives the upper bound in Theo-rem 1.2. To obtain the upper bound in Theorem 1.1, let N be given, and assume that N ⩾2(2L0). We can write N = ηM2M + ηM−12M−1 + · · · + η12 + η0 for numbers (ηM, . . . , η0) ∈{0, 1}M+1, where we assume that M is chosen in such a way that ηM = 1; this is simply the binary representation of N.
For simplicity of writing we set NM+1 = 0 and Nµ = ηM2M + ηM−12M−1 + · · · + ηµ2µ, 0 ⩽µ ⩽M.
Then clearly we have (4.9) N X k=1 e2πinkα = M+1 X µ=1 ηµ−1 Nµ−1 X k=Nµ+1 e2πinkα = M+1 X µ=1 ηµ−1 2µ−1 X k=1 e2πinNµ+kα.
Note that for the “odious numbers” mk we have, for every k, that mk = ( 2k −1 if nk = 2k −2, 2k −2 if nk = 2k −1, where nk are the “evil numbers”. Furthermore, from the special structure of the Thue–Morse sequence we see that for 1 ⩽k ⩽2µ−1 we have (4.10) nNµ+k = ( 2Nµ + nk if s2(Nµ) = 0, 2Nµ + mk if s2(Nµ) = 1, which together with (1.16), (1.18), (2.2) and (2.3) implies that 2µ−1 X k=1 e2πinNµ+kα ⩽ exp µ−1 X ℓ=0 f1(2ℓα) !
+ exp µ−1 X ℓ=0 f2(2ℓα) !!
, 1 ⩽µ ⩽M + 1. Thus by (4.9) and (4.10) we have N X k=1 e2πinkα ⩽ M+1 X µ=1 2µ−1 X k=1 e2πinNµ+kα (4.11) ⩽4M + M+1 X µ=⌈log2 M⌉ 2µ−1 X k=1 e2πinNµ+kα ⩽4M + M+1 X µ=⌈log2 M⌉ exp µ−1 X ℓ=0 f1(2ℓα) (4.12) + exp µ−1 X ℓ=0 f2(2ℓα) !
.
ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 665 Using the fact that the assumption N ⩾2(2L0) implies that log2 M ⩾L0, and also using the inequalities (4.7) and (4.8), we obtain (4.13) M+1 X µ=⌈log2 M⌉ exp µ−1 X ℓ=0 f1(2ℓα) !
⩽ M+1 X µ=⌈log2 M⌉ exp π √ 2 + ˆ ε p 2(µ −1) log log(µ −1) ⩽M exp π √ 2 + ˆ ε p 2M log log M and (4.14) M+1 X µ=⌈log2 M⌉ exp µ−1 X ℓ=0 f2(2ℓα) !
⩽ M+1 X µ=⌈log2 M⌉ exp ˆ ε p 2(µ −1) log log(µ −1) ⩽M exp ˆ ε p 2M log log M .
Combining (4.12), (4.13) and (4.14) we obtain N X k=1 e2πinkα ⩽4M + 2M exp π √ 2 + ˆ ε p 2M log log M ⩽4M + 2M exp π √ 2 + ˆ ε p 2(log2 N) log log(log2 N) .
As a consequence we have N X k=1 e2πinkα ⩽exp π √log 2 + ε p (log N) log log log N for all sufficiently large N. This proves the upper bound in Theorem 1.1.
Part 2: lower bounds. — Now we prove the lower bound in Theorems 1.1 and 1.2. Again we assume that h = 1, that α and ε > 0 are fixed, and that α is from the set of full measure for which the conclusion of Lemma 4.1 holds. Again we set ˆ ε = ε/2. Then by Lemma 4.1 there exist infinitely many values of L for which both inequalities L X ℓ=0 f1(2ℓα) ⩾ π √ 2 −ˆ ε p 2L log log L TOME 67 (2017), FASCICULE 2 666 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER and L X ℓ=0 f2(2ℓα) ⩽ˆ ε p 2L log log L hold simultaneously. The first relation already gives the lower bound in Theorem 1.2. By (1.16) and (2.2), (2.3) we have 2L X k=1 e2πinkα ⩾1 2 exp L X ℓ=0 f1(2ℓα) !
−1 2 exp L X ℓ=0 f2(2ℓα) !
.
Thus for infinitely many L we have (4.15) 2L X k=1 e2πinkα ⩾1 2 exp π √ 2 −ˆ ε p 2L log log L −1 2 exp ˆ ε p 2L log log L .
Consequently we also have N X k=1 e2πinkα ⩾exp π √log 2 −ε p (log N) log log log N for infinitely many N. This proves the lower bound in Theorem 1.1.
□ 5. Discrepancy of evil Kronecker sequences: Proof of Theorem 1.3 For given L ⩾1, we set I1(L) = Z 1 0 L−1 Y ℓ=0 2 sin(π2ℓα) !
dα and I2(L) = Z 1 0 L−1 Y ℓ=0 2 cos(π2ℓα) !
dα.
Integrals of this type have been studied in great detail in . For the integral I1 it is proved there that (5.1) 2−LI1(L) = κλL (1 + o(1)) where κ, λ are positive constants with 0.654336 ⩽λ ⩽0.663197 (thereby improving an earlier result of Èminyan ).
ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 667 Hence for every ε > 0 for L large enough we have (5.2) I1(L) ⩽ 2L1+log2 λ+ε .
We will improve the estimate given for λ by Fouvry and Mauduit in the following.
Lemma 5.1. — Let λ be defined as in (5.1). Then 0.66130 < λ < 0.66135.
Proof. — By the formula above of equation (4.2) in we have Z 1 0 L−1 Y ℓ=0 sin π2ℓα dα = Z 1 0 φj(α) L−j−1 Y ℓ=0 sin π2ℓα dα where φ0(α) ≡1 and φj+1(α) = 1 2 sin π α 2 φj α 2 + cos π α 2 φj α + 1 2 .
Furthermore, it was shown in that the functions φj are symmetric around α = 1 2 on [0, 1], and that they are concave on [0, 1]. Hence φj(0) = min α∈[0,1]φj(α), and φj (1/2) = max α∈[0,1]φj(α).
Let qj(α) := φj(α) φj−1(α), mj := min α∈[0,1]qj(α), and Mj := max α∈[0,1]qj(α). Note that qj(α) of course also is symmetric around α = 1 2 in [0, 1]. We have for every α ∈[0, 1] qj+1(α) = φj+1(α) φj(α) = sin π α 2 φj α 2 + cos π α 2 φj α+1 2 sin π α 2 φj−1 α 2 + cos π α 2 φj−1 α+1 2 ⩽ sin π α 2 Mjφj−1 α 2 + cos π α 2 Mjφj−1 α+1 2 sin π α 2 φj−1 α 2 + cos π α 2 φj−1 α+1 2 = Mj.
Therefore Mj+1 = max α∈[0,1]qj(α) ⩽Mj.
Analogously we obtain qj+1(α) ⩾mj for all α ∈[0, 1]. Altogether M1 ⩾ M2 ⩾M3 ⩾. . . and m1 ⩽m2 ⩽m3 ⩽. . ., and therefore for every k fixed TOME 67 (2017), FASCICULE 2 668 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER we have Z 1 0 L−1 Y ℓ=0 sin π2ℓα dα = Z 1 0 φL(α)dα = Z 1 0 L Y j=1 qj(α)dα ⩽M L−k k Z 1 0 q1(α) . . . qk−1(α) dα.
Similarly we get Z 1 0 L−1 Y ℓ=0 sin π2ℓα dα ⩾mL−k k Z 1 0 q1(α) . . . qk−1(α) dα.
Hence mL−k k Z 1 0 q1(α) . . . qk−1(α) dα ⩽κλL (1 + o(1)) ⩽M L−k k Z 1 0 q1(α) . . . qk−1(α) dα and consequently (5.3) mk ⩽λ ⩽Mk for all k ⩾1.
By considering the function q6(α), in the following we will prove that m6 > 0.6613 and M6 < 0.66135.
0.2 0.4 0.6 0.8 1.0 0.66130 0.66131 0.66132 0.66133 0.66134 Figure 5.1. The function q6(α).
ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 669 Wherever q6 is differentiable we have |q6 ′(α)| = φ6 ′(α)φ5(α) −φ6(α)φ5 ′(α) (φ5(α))2 ⩽ max α∈[0,1] φ6 ′(α) φ5(0) + φ6 1 2 (φ5(0))2 max α∈[0,1] φ5 ′(α) .
(5.4) It can easily be checked for example by differentiating φ5(α) and φ6(α) with the help of Mathematica that φ5 ′(α) is the sum resp. difference of – 32 products of absolute values of sines and cosines, each product weighted by a factor π 1024, – 32 such products weighted by π 512, – 32 with weight π 256, – 32 with weight π 128, – 32 products with factor π 64.
Hence φ5 ′(α) ⩽32π 1 1024 + 1 512 + 1 256 + 1 128 + 1 64 = 31 32π.
In the same way we show that also φ6 ′(α) ⩽31 32π. By combining these estimates with the values of φ5(0) and φ6(1/2) in (5.4), we finally obtain |q6′(α)| ⩽56.4.
Now we calculate q6 a 2800000 for a = 0, 1, . . . , 1400000 with the help of Mathematica and obtain max q6 a 2800000 a = 0, 1, . . . , 1400000 = 0.66133092 . . .
min q6 a 2800000 a = 0, 1, . . . , 1400000 = 0.66131148 . . .
Hence m6 ⩾0.66131145 −56.4 1 5600000 = 0.661301 . . .
and M6 ⩽0.66133092 + 56.4 1 5600000 = 0.661341 . . .
By (5.3) this implies Lemma 5.1.
□ Let us remark that numerical experiments with q15(α) suggest that λ = 0.661322602 . . . . It is tempting to conjecture that the precise value of λ can be expressed in a simple way in terms of the “usual” mathematical constants such as e, π, log 2, etc. However, we do not know what such an TOME 67 (2017), FASCICULE 2 670 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER expression could look like, and cannot even make a reasonable guess (the numerical argument in the proof of Lemma 5.1 does not give any hints).
Proof of Theorem 1.3.
Part 1: upper bound. — For the integral I2(L) we can use the equality (5.5) L−1 Y ℓ=0 2 cos(π2ℓα) = sin π2Lα sin πα ⩽min 2L, 1 πα(1 −α) which holds for 0 < α < 1 and which implies that I2(L) ⩽ Z 1 0 sin π2Lα sin πα dα ⩽2 Z 2−L 0 2L dα + Z 1−2−L 2−L 1 πα(1 −α) dα ⩽2 + 2L log 2 π (5.6) (this is essentially a variant of the classical bound for the L1-norm of the Dirichlet kernel).
For µ ⩾1 and ε > 0 we set (5.7) Gµ = α ∈(0, 1) : 24µ X h=1 1 h 2µ−1 X k=1 e2πihnkα > (2µ)1+log2 λ+ε .
By (1.16), (5.2) and (5.6) we have Z 1 0 24µ X h=1 1 h 2µ−1 X k=1 e2πihnkα dα = 24µ X h=1 1 h Z 1 0 2µ−1 X k=1 e2πihnkα dα ⩽ 24µ X h=1 1 h (I1(µ) + I2(µ)) ⩽(2µ)1+log2 λ+ ε 2 for sufficiently large µ. Consequently we have P(Gµ) ⩽(2µ)−ε 2 .
which implies that by the Borel–Cantelli lemma with probability one only finitely many events Gµ occur. We can show the same result if we replace the sequence (nk)k⩾1 in (5.7) by (mk)k⩾1. In other words, for almost all α we have (5.8) 24µ X h=1 1 h 2µ−1 X k=1 e2πihnkα ⩽(2µ)1+log2 λ+ε ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 671 and the same estimate for (mk)k⩾1 instead of (nk)k⩾1, for all µ ⩾µ0(α, ε).
Now assume that α, ε and N are given. Furthermore we assume that for these values of α and ε the estimate (5.8) and the corresponding estimate for (mk)k⩾1 instead of (nk)k⩾1 hold for µ ⩾µ0, and that N ⩾2(4µ0). We apply the same dyadic decomposition of N as in the proof of the upper bound of Theorem 1.1 in Section 4. In the same way as we obtained (4.11), together with the Koksma–Erdős–Turán inequality, we can now obtain with M = ⌊log2(N)⌋ (5.9) ND∗ N({n1α}, . . . , {nNα}) ⩽1 + N X h=1 1 h N X k=1 e2πihnkα ⩽1 + N X h=1 1 h M+1 X µ=1 2µ−1 X k=1 e2πihnNµ+kα ⩽1 + ⌈(log2 N)/4⌉ X µ=1 N X h=1 1 h 2µ−1 X k=1 e2πihnNµ+kα | {z } ≪N1/4 log N + M+1 X µ=⌈(log2 N)/4⌉+1 N X h=1 1 h 2µ−1 X k=1 e2πihnNµ+kα .
For the last term in (5.9) by (4.10) and (5.8) we have M+1 X µ=⌈(log2 N)/4⌉+1 N X h=1 1 h 2µ−1 X k=1 e2πihnNµ+kα ⩽ M+1 X µ=⌈(log2 N)/4⌉+1 24µ X h=1 1 h 2µ−1 X k=1 e2πihnNµ+kα ⩽ M+1 X µ=⌈(log2 N)/4⌉+1 (2µ)1+log2 λ+ε/2 | {z } ≪(2M)1+log2 λ+ε/2 .
Since M ⩽log2 N, together with (5.9) we have proved that ND∗ N({n1α}, . . . , {nNα}) ⩽N 1+log2 λ+ε for all sufficiently large N, which proves the upper bound in Theorem 1.3.
TOME 67 (2017), FASCICULE 2 672 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER Part 2: lower bound. — By (1.3) and (1.17) for the discrepancy D∗ N of the sequence ({nkα})k⩾1 with N = 2L for each positive integer H we have D∗ N ⩾ 1 4H 1 N N X k=1 e2πinkHα ⩾ 1 4H L Y ℓ=0 sin πH2ℓα −1 4H L Y ℓ=0 cos πH2ℓα | {z } ≪(H∥Hα∥)−1 .
Let fL(α) = L Y ℓ=0 sin π2ℓα .
We will show below that for any given ε > 0 for almost all α there are infinitely many L such that there exists a positive integer hL with hL ⩽2L and (5.10) 1 hL fL (hLα) ≫λL(1+ε).
It is a well-known fact in metric Diophantine approximation that for almost all α we have h ∥hα∥⩾ 1 hε for all h large enough. Hence if (5.10) is true for almost all α then there are infinitely many L such that for N = 2L we have D∗ N ≫λL(1+ε) −(hL)ε (5.11) ≫ 2L(1+ε) log λ log 2 − 2Lε ≫N (1+ε) log λ log 2 , and the desired result follows (note that log λ < 0). It remains to show the existence of the numbers hL ⩽2L which satisfy (5.10).
Let ε > 0 be given. From the definition of fL(α) it is easily seen that (5.12) |fL(α1) −fL(α2)| ⩽2L+1π|α1 −α2|; this follows from the fact that the derivative of the function QL ℓ=0 sin π2ℓα is bounded uniformly by 2L+1π. Now let gL(α) be the function defined by gL(α) = fL(j4−L) for α ∈ h j4−L, (j + 1)4−L , for j = 0, . . . , 4L −1.
This definition means that gL is constant on intervals of length 4−L which lie between two integer multiples of 4−L, and coincides with fL on the left endpoint of such intervals. By (5.12) we have (5.13) |gL −fL| ⩽2π2−L, ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 673 which means that it is sufficient to prove (5.10) with fL replaced by gL (remember that the value of ε > 0 was arbitrary and λ > 1/2). The reason for using the functions gL instead of fL is that every function gL can be written as a sum of at most 4L different characteristic functions of intervals; consequently, we know that the set of values of α where |gL| is “large” can be written as the union of at most 4L intervals, which implies an upper bound for the size of the Fourier coefficients of the characteristic function of this set (see below for details).
Let Q = Q (ε) be a positive integer which will be chosen in dependence on ε (we assume that Q is “large”). We define real numbers δi = λ1−i−1/2 Q , i = 0, 1, . . . , Q + 1.
Note that δ0 < δ1 < . . . < δQ+1, δ0 = λ1+1/(2Q) < λ, and δQ+1 > 1.
Furthermore, we define M (i) L := n α ∈[0, 1) : δL i < |gL (α) | ⩽δL i+1 o for i = 0, 1, . . . , Q.
Then by (5.1) and (5.13) we have Q X i=0 P M (i) L δL i+1 + 1 − Q X i=0 P M (i) L ! λ1+ 1 2Q L ⩾ Z 1 0 |gL (α) | dα > κ 2 λL for sufficiently large L, where κ is the number from (5.1). Hence we have Q X i=0 P M (i) L δL i+1 ⩾κ 4 λL for sufficiently large L. Consequently for every L large enough there is an iL ∈{0, . . . , Q} with δL iL+1P M (iL) L ⩾ κ 4(Q + 1)λL, which implies that P M (iL) L ⩾ κ 4(Q + 1) λ δiL+1 L .
Note that, as a consequence of the construction of gL, the set M (i) L always is a union of at most 4L disjoint intervals. It is easily seen that by trimming the sets M (i) L appropriately we can always find a set R(i) L such that R(i) L ⊂M (i) L , TOME 67 (2017), FASCICULE 2 674 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER such that R(i) L is also the union of at most 4L intervals, and such that for the measure of the sets R(i) L we have the exact equality P R(iL) L = κ 4(Q + 1) λ δiL+1 L .
Let η = η(ε) > 0 be a “small” number. Let HL denote the largest integer such that (5.14) HL ⩽4(Q + 1) κ δiL+1 λ L (1 + η)L.
Note that δiL+1 λ ⩾λ 1 2Q , and consequently we have HL ⩾(1 + η)L for L large enough; this means that HL grows exponentially in L. Note also that it is easily seen that HL ⩽2L for sufficiently large L (provided that η is chosen sufficiently small), which is important for (5.11).
We will show that for almost all α for infinitely many L there is a (5.15) h ⩽HL such that {hα} ∈R(iL) L .
For these h then we have |gL(hα)| h ⩾δL iL κ 4(Q + 1) λ δiL+1 L (1 + η)−L = κ 4(Q + 1) δiL δiL+1 L (1 + η)−LλL ≫λL(1+ε), for Q large enough and η small enough in dependence on ε. Together with (5.13) this will establish (5.10), as desired.
It remains to show (5.15). Let 1L(α) denote the characteristic function of the set R(i) L , extended with period one. Then we know that (5.16) Z 1 0 1L(α) dα = κ 4Q λ δiL+1 L .
Setting IL(α) = 1L(α) − Z 1 0 1L(ω) dω, we clearly have R 1 0 IL(α) dα = 0 and (5.17) Var[0,1] IL ⩽4L.
ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 675 From (5.16) we can easily calculate that ∥IL∥2 2 = Z 1 0 IL(α)2dα = κ 4Q λ δiL+1 L 1 −κ 4Q λ δiL+1 L!
⩽κ 4Q λ δiL+1 L .
(5.18) We write IL(α) ∼ ∞ X j=1 (aj cos 2πjα + bj sin 2πjα) for the Fourier series of IL (note that it has no constant term, since IL has integral zero). In the sequel, we want to show that the sum (5.19) X h⩽HL Z 1 0 1L(hα) dα is large in comparison with the sum (5.20) X h⩽HL IL(hα), for almost all α and infinitely many L. Since (5.21) X h⩽HL 1L(hα) = X h⩽HL Z 1 0 1L(ω) dω + X h⩽HL IL(hα), such an estimate will show that the sum on the left-hand side of (5.21) is large (for almost all α, for infinitely many L), which in turn implies that many of the events described in (5.15) will occur. A lower bound for (5.19) is easy to obtain; to find an asymptotic upper bound for (5.20), we will calculate the L2 norm of these sums, and apply the Borel–Cantelli lemma.
From (5.16) we directly obtain (5.22) HL X h=1 Z 1 0 1L(α) dα = HL κ 4Q λ δiL+1 L ≫(1 + η)L.
Next we estimate HL X h=1 IL(h·) 2 , which is relatively difficult. As a consequence of (5.17) and a classical in-equality for the size of the Fourier coefficients of functions of bounded variation (see for example [41, p. 48]) we have (5.23) |aj| ⩽Var[0,1]IL 2j ⩽4L j , and similarly |bj| ⩽4L j .
TOME 67 (2017), FASCICULE 2 676 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER We split the function IL into an even and an odd part (that is, into a cosine-and a sine-series). In the sequel, we consider only the even part; the odd part can be treated in exactly the same way. Let pL(α) denote the 43L-th partial sum of the Fourier series of the even part of IL, and let rL(α) denote the remainder term. Then by Minkowski’s inequality we have (5.24) HL X h=1 I(even) L (h·) 2 ⩽ HL X h=1 pL(h·) 2 + HL X h=1 rL(h·) 2 .
Furthermore, (5.23), Minkowski’s inequality, and Parseval’s identity imply that HL X h=1 rL(h·) 2 ⩽HL∥rL∥2 ⩽HL v u u t ∞ X j=43L+1 42L j2 ⩽HL2−L ≪1.
(5.25) To estimate the first term on the right-hand side of (5.24), we expand pL into a Fourier series and use the orthogonality of the trigonometric system.
Then we obtain (5.26) HL X h=1 pL(h·) 2 2 = HL X n1,n2=1 43L X j1,j2=1 | {z } j1n1=j2n2 aj1aj2 2 = 43L X j1,j2=1 aj1aj2 2 # n (n1, n2) : 1 ⩽n1, n2 ⩽HL, j1n1 = j2n2 o .
To estimate the size of the sum on the right-hand size of (5.26), we assume that j1 and j2 are fixed. In the case j1 = 1 and j2 = 1, we clearly have j1n1 = j2n2 whenever n1 = n2; thus the cardinality of the set on the right-hand side of (5.26) is HL. If j1 = 1 and j2 = 2, then we have to count the number of pairs (n1, n2) for which 2n1 = n2; this number is ⌊HL/2⌋.
For the values j1 = 2 and j2 = 4 we also have to count the number of pairs (n1, n2) for which 2n1 = n2; so this cardinality also is ⌊HL/2⌋. The ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 677 last example shows that the greatest common divisor of j1 and j2 plays a role in this calculation. Using similar considerations, in the case of general (fixed) values of j1 and j2 it turns out that we have j1n1 = j2n2 whenever n1 = v j2 gcd(j1, j2), n2 = v j1 gcd(v1, v2) for some positive integer v.
As a consequence we have # n (n1, n2) : 1 ⩽n1, n2 ⩽HL, j1n1 = j2n2 o = # v ⩾1 : v ⩽min HL gcd(j1, j2) j2 , HL gcd(j1, j2) j1 = HL gcd(j1, j2) max(j1, j2) ⩽HL gcd(j1, j2) √j1j2 .
Combining this estimate with (5.26) we obtain (5.27) HL X h=1 pL(h·) 2 2 ⩽HL 43L X j1,j2=1 |aj1aj2| 2 gcd(j1, j2) √j1j2 .
The sum on the right-hand side of the last equation is called a GCD sum.
It is well-known that such sums play an important role in the metric the-ory of Diophantine approximation; the particular sum in (5.27) probably appeared for the first time in LeVeque’s paper (see also and ). A precise upper bound for these sums has been obtained by Hilberdink .(3) Hilberdink’s result implies that there exists an absolute constant cabs such that 43L X j1,j2=1 |aj1aj2| 2 gcd(j1, j2) √j1j2 ≪exp cabs p log(43L) p log log 43L ! 43L X j=1 a2 j.
(3) The upper bounds in are formulated in terms bounds for the maximal size of the largest eigenvalue of certain matrices involving greatest common divisors (GCD matrices). In it is explained in detail how these bounds for eigenvalues of GCD matrices can be translated into bounds for GCD sums such as the ones considered here.
TOME 67 (2017), FASCICULE 2 678 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER Combining this estimate with (5.18) and (5.27) (and using Parseval’s iden-tity) we have HL X h=1 pL(h·) 2 2 ≪HL exp cabs p log(43L) p log log 43L !
κ 4Q λ δiL+1 L ≪(1 + η)L exp cabs p log(43L) p log log 43L !
, and, together with (5.24) and (5.25), and with a similar argument for the odd part of IL, we obtain (5.28) HL X h=1 IL(h·) 2 2 ≪(1 + η)L exp cabs p log(43L) p log log 43L !
.
By Chebyshev’s inequality we have P α ∈[0, 1) : HL X h=1 IL(hα) > (log HL) HL X h=1 IL(h·) 2 !
⩽ 1 (log HL)2 , and since (HL)L⩾1 grows exponentially in L these probabilities give a con-vergent series when summing over L. Thus by the Borel–Cantelli lemma with probability one only finitely many events HL X h=1 IL(hα) > (log HL) M X h=1 IL(h·) 2 happen, which by (5.28) implies that HL X h=1 IL(hα) ≪(1 + η)L/2 exp ˆ cabs √ L √log L !
for some absolute constant ˆ cabs. Comparing this upper bound with (5.22) and using (5.21) we conclude that HL X h=1 1L(hα) ≫(1 + η)L as L →∞ for almost all α. In particular we have ∞ X L=1 HL X h=1 1L(hα) = ∞ for almost all α, which means that for almost all α infinitely many events (5.15) occur. As noted after equation (5.15), this proves the the-orem.
□ ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 679 6. Concrete Examples: Proof of Theorem 1.4 It is known that for all α we have (6.1) L Y ℓ=0 sin π2ℓα ⩽HL for all L, where H = √ 3 2 = 0.866 . . .. This inequality appears for example as formula (2.10) in , but was already proved earlier in and .
Thus from (1.16), (5.5) and the Weyl criterion it follows that ({nkα})k⩾1 is u.d. mod 1 iffα is irrational.
Hence by (1.1), (1.3), (1.16), and (1.17) for N of the form N = 2L we have max h⩽N 1 h L Y l=0 |2 sin πh2lα| − L Y l=0 |2 cos πh2lα| ≪N e D∗ N ≪ N X h=1 1 h ∥hα∥+ N X h=1 1 h L Y k=0 2 sin π2khα ≪ N X h=1 1 h ∥hα∥+ (log N)N log 3 log 4 , where we write e D∗ N for the star-discrepancy of the first N terms of the evil Kronecker sequence.
From the left-hand side of this inequality it is not difficult – but we do not want to go into the details here – to show that max n⩽N n e D∗ n ≫ 1 log N max n⩽N nD∗ n for all N, where D∗ N denotes the star-discrepancy of the pure Kronecker sequence.
Furthermore it is easy to show – we again do not go into the details – that N X h=1 1 h ∥hα∥≪(log N)ND∗ N for all N = 2L.
Hence for all N we have (6.2) 1 log N max n⩽N nD∗ n ≪max n⩽N n e D∗ n ≪(log N)2 max n⩽N nD∗ n+(log N)2 N log 3 log 4 .
From this we conclude that the order of the discrepancy N e D∗ N of the evil Kronecker sequence always is essentially (up to logarithmic factors) larger TOME 67 (2017), FASCICULE 2 680 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER or equal to the order of the discrepancy ND∗ N of the pure Kronecker se-quence. In fact, the order of N e D∗ N essentially equals the order of ND∗ N plus an expression which is at most of order N log 3 log 4 +ε. The order of the additional expression is controlled by lacunary products of sine-functions.
Hence we conclude what we have already announced in Section 1: – If the order of D∗ N satisfies ND∗ N = Ω N log 3 log 4 then e D∗ N essentially is of the same order as D∗ N – If D∗ N satisfies ND∗ N = O N log 3 log 4 then e D∗ N satisfies N e D∗ N = O N log 3 log 4 +ε Hence the two examples given in Theorem 1.4 show interesting non-trivial cases where we have (almost) best possible distribution for the pure Kro-necker sequence with bad distribution for the evil Kronecker sequence. In-deed especially the first example gives essentially the extremal values for D∗ N and for e D∗ N, and shows that the right-hand side of (6.2) is also essen-tially optimal.
It remains an open problem to give concrete examples α where the cor-responding evil Kronecker sequence has “small” discrepancy D∗ N, e.g., a discrepancy of the metric order given in Theorem 1.3 or smaller. Of course it also remains an open problem to give good estimates for e D∗ N in the case of “natural” examples of α like α = √ 2.
Proof of Theorem 1.4.
(a). — By we know that the continued fraction coefficients of the number β := P∞ k=1 1 42k are bounded. This is equivalent to the existence of some c > 0 such that β −p q > c q2 for all p, q ∈Z, q ⩾1, i.e., |qβ −p| > c q for all such p and q.
This implies especially |3qβ + 2q −3p| > c 3q for all p, q ∈Z, q ⩾1, i.e., β + 2 3 −p q > c 9 q2 for all such p, q, and hence α = β + 2 3 has bounded con-tinued fraction coefficients. So the star-discrepancy of the pure Kronecker sequence ({nα})n⩾1 satisfies ND∗ N = O (log N).
On the other hand we already know that for the star-discrepancy e D∗ N of the evil Kronecker sequence ({nkα})k⩾1 with N = 2L we have N e D∗ N ≫ N X k=1 exp (2πinkα) ≫ L Y ℓ=0 2 sin π2ℓα − 1 ∥α∥ .
ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 681 We give a suitable lower bound for ΠL := L Y ℓ=0 2 sin π2ℓα = L Y ℓ=0 2 sin π 3 L Y ℓ=0 sin π2ℓα sin π 3 = √ 3N log 3 log 4 L Y ℓ=0 cos (πδℓ) + (−1)ℓ+1 √ 3 sin (πδℓ) .
Here δℓ:= 2ℓβ and we have used sin (x + y) = sin x · cos y + cos x · sin y.
Note that the base 2 representation of 2ℓα has one of the following ten possible forms: 0.1010 . . .
0.0101 . . .
0.0010 . . .
0.1101 . . .
0.1110 . . .
0.1001 . . .
0.1100 . . .
0.0111 . . .
0.1011 . . .
0.0110 . . .
hence 2ℓα > 1 16 always and therefore |sin π2ℓα| |sin π 3 | > 0.2 always.
Because of |cos πx −1| ⩽3x and |sin πx| ⩽πx for x ⩾0 we have cos πδℓ+ (−1)ℓ+1 √ 3 sin πδℓ ⩾1 − 3 + π √ 3 δℓ> 1 −5δℓ.
Therefore, noting that max(0.2, 1 −5x) > e−11x for x > 0, we also have L Y ℓ=0 sin π2ℓα sin π 3 > L Y ℓ=0 max (0.2, 1 −5δℓ) > e−11PL ℓ=0 δℓ ≫e−22 log L.
TOME 67 (2017), FASCICULE 2 682 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER So ΠL ≫ 1 L22 N log 3 log 4 ≫ 1 (log N)22 N log 3 log 4 and the lower bound for e D∗ N follows. The upper bound for e D∗ N follows from (6.2) and from the upper bound for D∗ N.
(b). — It was shown in that γ has approximation degree 1, hence for the star-discrepancy of the sequence ({nα})n⩾1 we have ND∗ N = O (N ε) for every ε > 0. To prove the lower bound for the star-discrepancy e D∗ N of the sequence ({nkγ})k⩾1, like in the proof of part (a) we have to estimate ΠL := QL ℓ=0 2 sin π2ℓγ from below. We will give in the following as an additional information also an upper estimate for ΠL in order to show that our lower estimate is rather sharp.
We may restrict ourselves to L of the form L = 8U −1. Then ΠL = U−1 Y j=0 8j+7 Y ℓ=8j 2 sin π2ℓγ .
In the following we use some well-known facts on properties of the Thue– Morse sequence: the base 2-representation of γ = 0, γ1γ2γ3 . . . consists of 8-blocks γ8v+1 . . . γ8v+8 of the form A := 10010110 or B := 01101001.
Four such consecutive 8-blocks can occur in the following ten combina-tions: c1 = AABA c2 = AABB c3 = ABAA c4 = ABBA c5 = ABAB c6 = BBAB c7 = BBAA c8 = BABB c9 = BAAB c10 = BABA Let, for example, j be such that 28j = 0, c1 . . . .
(c1 is the block of 32 digits defined above) ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 683 For m = 0, 1, 2, . . . , 7 let x1,m := ( 2m · 0, c1 if {2m · 0, c1} < 1 2 2m · 0, c1 + 1 232 if {2m · 0, c1} > 1 2 and y1,m := ( 2m · 0, c1 if {2m · 0, c1} > 1 2 2m · 0, c1 + 1 232 if {2m · 0, c1} < 1 2 Then 8j+7 Y ℓ=8j 2 sin π2ℓγ < 7 Y m=0 |2 sin πy1,m| =: U(c1) = 33.487710 . . .
8j+7 Y ℓ=8j 2 sin π2ℓγ > 7 Y m=0 |2 sin πx1,m| =: D(c1) = 33.487705 . . .
In the same way we determine U (ci) and d (ci) for i = 2, 3, . . . , 10.
Using the special structure of the Thue–Morse sequence as a sequence gen-erated by a finite automaton one can use the methods from to show that the asymptotic frequencies F (ci) of the occurrence of a quadruple ci of 8-blocks in the Thue–Morse sequence are given by F (c1) = F (c2) = F (c3) = F (c5) = F (c6) = F (c7) = F (c8) = F (c10) = 1 12 and F (c4) = F (c9) = 1 6.
Hence we get (1 −ϵ)U · (D (c1) D (c2) D (c3) D (c5) D (c6) D (c7) D (c8) D (c10)) U 12 · (D (c4) D (c9)) U 6 ≪ L Y ℓ=0 2 sin π2ℓγ ≪(1 + ϵ)U · (U (c1) U (c2) U (c3) U (c5) U (c6) U (c7) U (c8) U (c10)) U 12 · (U (c4) U (c9)) U 6 TOME 67 (2017), FASCICULE 2 684 Christoph AISTLEITNER, Roswitha HOFER & Gerhard LARCHER which leads to N 0.6178775 ≪ L Y ℓ=0 2 sin π2ℓγ ≪N 0.6178777.
for ϵ small and L large enough. This finishes the proof.
□ 7. An open problem from the theory of metric Diophantine approximation In conclusion, we mention an open problem from the theory of Diophan-tine approximation which is related to our proof of the lower bound in Theorem 1.3. In metric Diophantine approximation, one is often interested in finding conditions on (φ(q))q⩾1 which guarantee that α −p q < φ(q) q has infinitely many integer solutions p, q for almost all α. Two instances of this problem, either under the additional requirement that p, q are co-prime (Duffin–Schaeffer conjecture) or without this additional requirement (Catlin conjecture), constitute probably the two most important open prob-lems in metric number theory. For the origin of the Duffin–Schaeffer con-jecture see , for the Catlin conjecture see . Problems of this type are discussed in great detail in Glyn Harman’s monograph on Metric Number Theory . For a recent survey, see .
The problem without the requirement of coprime solutions can also be written in the following form: Let A1, A2, . . . be intervals of length ⩽1, which are symmetric around 0. Let ψ1, ψ2, . . . denote the Lebesgue measure (that is, the length) of these intervals. Under which conditions on ψ1, ψ2, . . .
do we have ∞ X n=1 1An(nα) = ∞ for almost all α? Here 1A denotes the characteristic function of A, extended with period one.
Now in a first step this problem can be generalized to the case when the intervals A1, A2, . . . are not necessarily symmetric around 0, which leads to a problem in inhomogeneous Diophantine approximation. This type of question is also quite well-investigated.
Perpetuating this line of thought, it is natural to ask what happens if we do not assume that A1, A2, . . . are intervals, but if they may denote any ANNALES DE L’INSTITUT FOURIER EVIL KRONECKER SEQUENCES 685 measurable sets in [0, 1]. Writing ψ1, ψ2, . . . for the measure of these sets, the question is under which conditions on ψ1, ψ2, . . . we have ∞ X n=1 1An(nα) = ∞ for almost all α. Note that a necessary condition is the divergence of the sum of the measures, by the Borel–Cantelli lemma. It seems that hardly anything is known about this general problem. As far as we know, this problem was first stated by LeVeque in . In this paper he answered a conjecture of Erdős, and he formulated a generalized version of Erdős’ conjecture. We consider this as a very interesting open problem, and we re-state it below.
Open problem. — Let A1, A2, . . . be measurable sets in [0, 1], and let ψ1, ψ2, . . . denote their measure. Under which conditions on (ψn)n⩾1 is it true that for almost all α the fractional part {nα} is contained in the set An for infinitely many indices n; equivalently, under which conditions is it true that ∞ X n=1 1An(nα) = ∞ almost everywhere, where the characteristic functions are extended with period one.
A problem quite similar to this one emerged during the proof of the lower bound of Theorem 1.3. However, the situation was comparatively simple there, for example since we could assume there that the sets An can be written as the sum of a moderate number of intervals. The general problem seems to be much more complicated; LeVeque wrote that this general problem “seems rather intractable”.
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Manuscrit reçu le 11 janvier 2016, révisé le 7 avril 2016, accepté le 12 mai 2016.
Christoph AISTLEITNER Institute of Financial Mathematics and Applied Number Theory Johannes Kepler University Linz Altenbergerstr. 69, 4040 Linz (Austria) aistleitner@math.tugraz.at Roswitha HOFER Institute of Financial Mathematics and Applied Number Theory Johannes Kepler University Linz Altenbergerstr. 69, 4040 Linz (Austria) roswitha.hofer@jku.at Gerhard LARCHER Institute of Financial Mathematics and Applied Number Theory Johannes Kepler University Linz Altenbergerstr. 69, 4040 Linz (Austria) gerhard.larcher@jku.at TOME 67 (2017), FASCICULE 2 |
13726 | https://math.stackexchange.com/questions/3847560/how-can-i-mathematiically-enforce-a-set-of-constraints-on-a-range-from-0-to-1 | How can I mathematiically enforce a set of constraints on a range from 0 to 1? - Mathematics Stack Exchange
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How can I mathematiically enforce a set of constraints on a range from 0 to 1?
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I have a function f(x)f(x) which needs to be bounded between 2 functions g(x)g(x) and h(x)h(x). Functions g(x)g(x) and h(x)h(x) are guaranteed never to intersect.
I have a function f(x)=a+b x+c x 2 f(x)=a+b x+c x 2 and I need to select parameters a a, b b, and c c such that ∫1 0 f(x)d x=w∫0 1 f(x)d x=w and g(x)<f(x)<h(x)g(x)<f(x)<h(x) for x∈[0,1]x∈[0,1].
How would I construct f(x) to enforce these constraints? All functions here are polynomials.
Edit: I have a hypothesis of how I might solve this, but I don't know how to formalize it:
If I specify this as an optimization problem,
min a,b,c(∫1 0 f(x)d x−w)2 min a,b,c(∫0 1 f(x)d x−w)2
and apply g(x) and h(x) as constraints to that optimization, perhaps such that I integrate f(x)−h(x)f(x)−h(x) only if the constraint is broken (and the same for g(x)g(x)) I could do something with Lagrange multipliers and slack variables... but I'm not sure how I would apply that.
constraints
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edited Oct 1, 2020 at 17:51
Anon YmousAnon Ymous
asked Oct 1, 2020 at 14:18
Anon YmousAnon Ymous
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Can you give the explicit formulas for g g and h h? It is not solvable in general.Federico –Federico 2020-10-01 15:43:38 +00:00 Commented Oct 1, 2020 at 15:43
The function g(x)g(x) is a parabola similar to x 2 x 2, the function h(x)h(x) is a vertical line above h(0)=3 h(0)=3. Perhaps my constraints are not well written, as the constraint is simply that f(x) does not intersect h(x) anywhere.Anon Ymous –Anon Ymous 2020-10-01 15:52:39 +00:00 Commented Oct 1, 2020 at 15:52
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When g g and h h form a narrow hose going up-down-up there will be no quadratic polynomial f f in between.
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answered Oct 1, 2020 at 15:37
Christian BlatterChristian Blatter
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Yes, but not on the interval x∈[0,1]x∈[0,1]Anon Ymous –Anon Ymous 2020-10-01 15:46:53 +00:00 Commented Oct 1, 2020 at 15:46
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Statistics
Introduction
StatisticsIntroduction
Contents Contents
Highlights
Table of contents
Preface
1 Sampling and Data
2 Descriptive Statistics
3 Probability Topics
4 Discrete Random Variables
5 Continuous Random Variables
6 The Normal Distribution
7 The Central Limit Theorem
8 Confidence Intervals
9 Hypothesis Testing with One Sample
10 Hypothesis Testing with Two Samples
11 The Chi-Square Distribution
12 Linear Regression and Correlation
Introduction
12.1 Linear Equations
12.2 The Regression Equation
12.3 Testing the Significance of the Correlation Coefficient (Optional)
12.4 Prediction (Optional)
12.5 Outliers
12.6 Regression (Distance from School) (Optional)
12.7 Regression (Textbook Cost) (Optional)
12.8 Regression (Fuel Efficiency) (Optional)
Key Terms
Chapter Review
Formula Review
Practice
Homework
Bringing It Together: Homework
References
Solutions
13 F Distribution and One-way Anova
A | Appendix A Review Exercises (Ch 3–13)
B | Appendix B Practice Tests (1–4) and Final Exams
C | Data Sets
D | Group and Partner Projects
E | Solution Sheets
F | Mathematical Phrases, Symbols, and Formulas
G | Notes for the TI-83, 83+, 84, 84+ Calculators
H | Tables
Index
Search for key terms or text.
Close
Figure 12.1 Linear regression and correlation can help you determine whether an auto mechanic’s salary is related to his work experience. (credit: Joshua Rothhaas)
Chapter Objectives
By the end of this chapter, the student should be able to do the following:
Discuss basic ideas of linear regression and correlation
Create and interpret a line of best fit
Calculate and interpret the correlation coefficient
Calculate and interpret outliers
Professionals often want to know how two or more numeric variables are related. For example, is there a relationship between the grade on the second math exam a student takes and the grade on the final exam? If there is a relationship, what is the relationship, and how strong is it?
In another example, your income may be determined by your education, your profession, your years of experience, and your ability. The amount you pay a repair person for labor is often determined by an initial amount plus an hourly fee.
The type of data described in the examples is bivariate data—bi—for two variables. In reality, statisticians use multivariate data, meaning many variables.
In this chapter, you will study the simplest form of regression—linear regression—with one independent variable (x). This involves data that fit a line in two dimensions. You will also study correlation, which measures the strength of a relationship.
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13728 | https://chem.libretexts.org/Courses/Barstow_Community_College/Survey_of_Chemistry_and_Physics/04%3A_Motion_Forces_and_Energy/4.01%3A_Motion/4.1.01%3A_Position_Displacement_and_Distance | Skip to main content
4.1.1: Position, Displacement, and Distance
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: Jun 6, 2025
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4.1: Motion
4.1.2: Time, Velocity, and Speed
Page ID
: 494311
OpenStax
OpenStax
( \newcommand{\kernel}{\mathrm{null}\,})
Learning Objectives
Define position, displacement, distance, and distance traveled.
Explain the relationship between position and displacement.
Distinguish between displacement and distance traveled.
Calculate displacement and distance given initial position, final position, and the path between the two.
Position
In order to describe the motion of an object, you must first be able to describe its position—where it is at any particular time. More precisely, you need to specify its position relative to a convenient reference frame. Earth is often used as a reference frame, and we often describe the position of an object as it relates to stationary objects in that reference frame. For example, a rocket launch would be described in terms of the position of the rocket with respect to the Earth as a whole, while a professor’s position could be described in terms of where she is in relation to the nearby white board. (See Figure 4.1.1.2.) In other cases, we use reference frames that are not stationary but are in motion relative to the Earth. To describe the position of a person in an airplane, for example, we use the airplane, not the Earth, as the reference frame. (See Figure4.1.1.3.)
Displacement
If an object moves relative to a reference frame (for example, if a professor moves to the right relative to a white board or a passenger moves toward the rear of an airplane), then the object’s position changes. This change in position is known as displacement. The word “displacement” implies that an object has moved, or has been displaced.
Definition: DISPLACEMENT
Displacement is the change in position of an object:
Δ𝑥=𝑥f−𝑥0,
where Δ𝑥 is displacement, 𝑥f is the final position, and 𝑥0 is the initial position.
In this text the upper case Greek letter Δ always means “change in” whatever quantity follows it; thus, Δ𝑥 means change in position. Always solve for displacement by subtracting initial position 𝑥0 from final position 𝑥𝑓.
Note that the SI unit for displacement is the meter (m), but sometimes kilometers, miles, feet, and other units of length are used. Keep in mind that when units other than the meter are used in a problem, you may need to convert them into meters to complete the calculation.
Note that displacement has a direction as well as a magnitude. The professor’s displacement is 2.0 m to the right, and the airline passenger’s displacement is 4.0 m toward the rear. In one-dimensional motion, direction can be specified with a plus or minus sign. When you begin a problem, you should select which direction is positive (usually that will be to the right or up, but you are free to select positive as being any direction). The professor’s initial position is 𝑥0 =1.5 m and her final position is 𝑥f =3.5 m. Thus her displacement is
Δ𝑥=𝑥f−𝑥0=3.5 m−1.5 m=+2.0 m.
In this coordinate system, motion to the right is positive, whereas motion to the left is negative. Similarly, the airplane passenger’s initial position is 𝑥0 =6.0 m and his final position is 𝑥f =2.0 m, so his displacement is
Δ𝑥=𝑥f−𝑥0=2.0 m−6.0 m=−4.0 m.
His displacement is negative because his motion is toward the rear of the plane, or in the negative 𝑥 direction in our coordinate system.
Distance
Although displacement is described in terms of direction, distance is not. Distance is defined to be the magnitude or size of displacement between two positions. Note that the distance between two positions is not the same as the distance traveled between them. Distance traveled is the total length of the path traveled between two positions. Distance has no direction and, thus, no sign. For example, the distance the professor walks is 2.0 m. The distance the airplane passenger walks is 4.0 m.
MISCONCEPTION ALERT: DISTANCE TRAVELED VS. MAGNITUDE OF DISPLACEMENT
It is important to note that the distance traveled, however, can be greater than the magnitude of the displacement (by magnitude, we mean just the size of the displacement without regard to its direction; that is, just a number with a unit). For example, the professor could pace back and forth many times, perhaps walking a distance of 150 m during a lecture, yet still end up only 2.0 m to the right of her starting point. In this case her displacement would be +2.0 m, the magnitude of her displacement would be 2.0 m, but the distance she traveled would be 150 m. In kinematics we nearly always deal with displacement and magnitude of displacement, and almost never with distance traveled. One way to think about this is to assume you marked the start of the motion and the end of the motion. The displacement is simply the difference in the position of the two marks and is independent of the path taken in traveling between the two marks. The distance traveled, however, is the total length of the path taken between the two marks.
Exercise 4.1.1.1
A cyclist rides 3 km west and then turns around and rides 2 km east. (a) What is her displacement? (b) What distance does she ride? (c) What is the magnitude of her displacement?
Answer
: (a) The rider’s displacement is Δ𝑥 =𝑥f −𝑥0 =−1 km. (The displacement is negative because we take east to be positive and west to be negative.)
: (b) The distance traveled is 3 km + 2 km = 5 km.
: (c) The magnitude of the displacement is 1 km.
Section Summary
Kinematics is the study of motion without considering its causes. In this chapter, it is limited to motion along a straight line, called one-dimensional motion.
Displacement is the change in position of an object.
In symbols, displacement Δx is defined to be
Δ𝑥=𝑥f−𝑥0,
where 𝑥0 is the initial position and 𝑥𝑓 is the final position. In this text, the Greek letter Δ (delta) always means “change in” whatever quantity follows it. The SI unit for displacement is the meter (m). Displacement has a direction as well as a magnitude.
When you start a problem, assign which direction will be positive.
Distance is the magnitude of displacement between two positions.
Distance traveled is the total length of the path traveled between two positions.
Glossary
kinematics
: the study of motion without considering its causes
position
: the location of an object at a particular time
displacement
: the change in position of an object
distance
: the magnitude of displacement between two positions
distance traveled
: the total length of the path traveled between two positions
4.1: Motion
4.1.2: Time, Velocity, and Speed |
13729 | https://study.com/academy/lesson/video/is-a-pentagon-a-regular-polygon.html | Pentagon Polygon | Definition, Types & Examples - Video | Study.com
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Video Summary for Regular Pentagon
A pentagon is a polygon with five straight sides and five angles, totaling 540 degrees.
The term comes from "penta" (five) and "gon" (angles).
Pentagons can be either regular or irregular.
Regular pentagons have equal sides and equal angles (108 degrees each).
Irregular pentagons have unequal sides and angles.
All polygons are closed shapes with straight sides that meet at points.
Regular polygons are both equiangular and equilateral
Irregular polygons have different measurements for sides and angles
This video demonstrates how to determine if a pentagon is regular by measuring its sides and angles.
When all sides and angles are equal, the pentagon is regular, otherwise it's irregular.
Read Pentagon Polygon | Definition, Types & Examples Lesson
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13730 | https://www.mi.sanu.ac.rs/vismath/wertheim/visualization3.HTM | vizualization vs. verbalization
2. How many convex Deltahedra possibly exist? An n-face Deltahedron is a polyhedron with n congruent equilateral triangular faces, conventionally denoted D n. (The name comes from the Greek capital letter Delta, which has a triangular shape.) An interesting problem is to find how many such Convex Polyhedra there are. In trying to exhaust the set of convex Deltahedra, we note, first of all, that this set includes three of the five Regular/Platonic solids: > > D 4 - the tetrahedron;D 8 - the octahedron;D 20 - the icosahedron Other convex Deltahedra might be obtainable by attaching several triangles side to side in different spatial configurations. Our question is: How many convex Deltahedra can thus be obtained? #### To be able to see VMRL images from this paper please download COSMO PLAYER 2.1 How many faces can convex Deltahedra possess? To answer this question we first observe that in each vertex 3, 4, or 5 triangular faces only, can meet. These numbers define the valence of a vertex. Clearly, the valences of all vertices in any given Deltahedron are not necessarily the same. Denote V 3,V 4, V 5-the number of vertices of valence 3, 4, 5 respectively. Next, note that (i) The minimum number of faces, is four, as in the Tetrahedron. (ii) Since no vertex can have more than five equilateral triangles meeting in it (as six give a flat 360 0 angle sum around the vertex), the largest possible convex Deltahedron is a solid having 5 equilateral triangles meet in each of its vertices, i.e., the Icosahedron. Thus twenty is the maximum number of faces for any Deltahedron. Therefore, The total number of faces of any existing convex Deltahedron is in the range between 4 and 20. 2.2 Candidates for existing Deltahedra First observe that the number of faces in any Deltahedron must be even. To establish it note that each face has three edges and each edge is shared by two faces, hence the relationship between the number of faces F and the number of edges E, satisfies:. Substituting this in Euler�s formula for simple Polyhedra: > F(aces)+V(ertices)=E(dges)+2 , we get > F=2(V-2) which implies the evenness of the total number of faces. Thus,we conclude that there are only 9 candidates for Deltahedra � those having 4, 6, 8, 10, 12, 14, 16, 18 and 20 (triangular) faces. To determine the various possible Deltahedra, we employ Descarte's formula for the total angular deficit of a convex polyhedron. First observe that: (i) The angular deficit of a vertex of valence 3 is 180 0 (ii) The angular deficit of a vertex of valence 4 is 120 0 (iii) The angular deficit of a vertex of valence 5 is 60 0 Descarte's formula for the total angular deficit of a convex polyhedron yields > 180 0 V 3+ 120 0 V 4 + 60 0 V 5 = 720 0 > > > or > > > 3V 3 + 2V 4 + V 5 = 12 Every Deltahedron must satisfy this (Diophantine) equation. This is a necessary but yet insufficient condition for the existence of a Deltahedron. 2.3 The solution Every solution of the above mentioned (Diophantine) equation, which takes into account the evenness of the total number of faces, is a candidate for an existing Deltahedron. Altogether, there are nineteen possible solutions to the equation. They are listed in Table 1 (Lichtenberg, 1988). However, the only solutions that correspond to actually existing Deltahedra are 1, 3-7,16, and 19, total of eight convex Deltahedra. > Visual and verbal description of the eight Deltahedra > > > NEXT > --------------------------------------------------------------------- > > > > Interestingly, if we do not restrict ourselves to convex deltahedra, there are infinitely many deltahedra. E.g., the stella octangula. The Angular Deficit or spherical deviation d - at a vertex of a polyhedron, defined as the difference between 360 0 and the sum of the angles of the polygons surrounding the vertex. The total angular deficit taken over each vertex of a convex polyhedron equals 720 0 |
13731 | https://artofproblemsolving.com/wiki/index.php/Incircle?srsltid=AfmBOopdZYTMwzxUdC420TvO1dcWs2HpAwSAuupsYDUeN8dS1xo02szJ | Art of Problem Solving
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Incircle
An incircle of a convexpolygon is a circle which is inside the figure and tangent to each side. Every triangle and regular polygon has a unique incircle, but in general polygons with 4 or more sides (such as non-squarerectangles) do not have an incircle. A quadrilateral that does have an incircle is called a Tangential Quadrilateral. For a triangle, the center of the incircle is the Incenter, where the incircle is the largest circle that can be inscribed in the polygon. The Incenter can be constructed by drawing the intersection of angle bisectors.
Formulas
The radius of an incircle of a triangle (the inradius) with sides and area is
The area of any triangle is where is the Semiperimeter of the triangle.
The formula above can be simplified with Heron's Formula, yielding
The radius of an incircle of a right triangle (the inradius) with legs and hypotenuse is .
For any polygon with an incircle, , where is the area, is the semi perimeter, and is the inradius.
The coordinates of the incenter (center of incircle) are , if the coordinates of each vertex are , , and , the side opposite of has length , the side opposite of has length , and the side opposite of has length .
The formula for the semiperimeter is .
The area of the triangle by Heron's Formula is .
See also
Circumradius
Inradius
Kimberling center
Circumcircle
Click here to learn about the orthocenter, and Line's Tangent
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13732 | http://www.cse.unl.edu/~goddard/Courses/CSCE310J/Lectures/Lecture2a-LogicIntro.pdf | 1 1 An Introduction to Logic By Chuck Cusack 2 Logic: Basic Definitions • Definition: A proposition is a statement that is either true or false, but not both.
• Defintion: The value of a proposition is called its truth value. Denoted by T if it is true, F if it is false Example 1: The statement “John Cusack is the president of the U.S.A.” is a proposition with truth value false.
Example 2: The statement “Do your homework” is not a proposition because it is not a statement that can be true or false.
3 Logical Connectives • Connectives are used to create a proposition from several other propositions.
• Such propositions are called compound propositions • The most common connectives are: – NEGATION (¬ or !) – AND (∧) – OR (∨) – XOR (⊕) – IMPLICATION (→) – BICONDITIONAL or IF AND ONLY IF (↔) 4 Connective Examples • Let p be the proposition “The sky is clear.” • Let q be the proposition “It is raining.” • Some examples that combine these are: – The sky is clear and it is raining. (p∧q) – The sky is clear and it is not raining. (p∧¬q) – It is raining if and only if the sky is not clear. (q ↔¬ p) 5 Truth Tables • Truth Tables are used to show the relationship between the truth values of individual propositions and the compound propositions based on them.
• Example: F T F F F F F F T T T T p ∧q q p 6 NEGATION • If p is a proposition, the negation of p, denoted ¬p, is “it is not the case that p.” • Example: Let p be the statement “this class has 30 students.” Then ¬p is the statement “this class does not have 30 students.” • It should be obvious that the negation of a proposition has the opposite truth value. In other words, if p is true, then ¬p is false.
• The truth table for ¬p is T F F T ¬p p 2 7 AND • Let p and q be propositions. The proposition “p and q,” denoted by p∧q, is true if and only if both p and q are true.
• p∧q is called the conjunction of p and q.
• The truth table for p∧q is F T F F F F F F T T T T p ∧q q p 8 OR • Let p and q be propositions. The proposition “p or q,” denoted by p∨q, is false if and only if both p and q are false. In other words, it is true if either p or q is true, and false otherwise.
• p∨q is called the disjunction of p and q.
• The truth table for p∨q is T T F F F F T F T T T T p∨q q p 9 XOR • Let p and q be propositions. The proposition “p exclusive or q,” denoted by p⊕q, is true if and only if either p or q is true, but not both.
• When the term OR is used in conversation, often the correct interpretation is XOR.
• The truth table for p⊕q is T T F F F F T F T F T T p⊕q q p 10 IMPLICATION • Let p and q be propositions. The proposition “p implies q,” denoted by p→q, is false if and only if p is true and q is false. • p→q is called an implication.
• The truth table for p→q is T T F T F F F F T T T T p→q q p 11 BICONDITIONAL • Let p and q be propositions. The proposition “p if and only if q,” denoted by p↔q, is true if and only if p and q have the same truth value. • p↔q is called a biconditional.
• The truth table for p↔q is F T F T F F F F T T T T p↔q q p 12 Constructing Truth Tables • Construct the truth table for the proposition ((p∧q)∨¬q) • We do this step by step as follows: T F F F F T T T q p F F F T p ∧q T F T F ¬q T F T T ((p∧q)∨¬q) 3 13 Everyday Logic • Logic is used in many places: – Writing – Speaking – Search engines – Mathematics – Computer Programs • A proper understanding of logic is useful, as the following examples will demonstrate.
14 Logic in Searching I • Situation: You want to find out all you can about disc golf.
• Problem: When you search for “disc golf,” you get many hits about golf and some about discs, but can’t find those about “disc golf.” • Solution: You need to find sites which mention both disc and golf, not either word. Search for disc AND golf 15 Logic in Searching II • Situation: You just bought some fresh corn, and you need a cornhusker to husk it, so you search for “cornhusker” on the Internet • Problem: Most of the results you get are about UNL’s football team.
• Solution: You need to find sites which mention cornhusker, but not UNL or football. Search for cornhusker AND NOT (UNL OR football) 16 Logic at Home • Situation: Your mom said “If you are good, you can have some ice cream or some cake.” • Problem: You were good, so you ate some ice cream and some cake. Your mom got mad because you had both.
• Solution: A simple miscommunication. By having ice cream and cake, you had ice cream or cake. But as is often the case in conversation, she really meant XOR, not OR. 17 Logic in School • Situation: You have 3 tests for a class. If you get an A on any two of them, or get an A on at least one but do not fail any of them, you will get an A for the course. • Problem: You are lazy, but want an A.
• Solution: Because of the OR condition, the minimal you can do is get an A on two exams and fail the third, or get an A on one exam and Ds on the other two. I’ll pick one A and 2 Ds.
18 Logic in Programming I • Situation: If x is greater than 0 and is less than or equal to 10, you need to increment it.
• Problem: You tried the following, but it seems too complicated, and doesn’t compile.
if(0<x<10 OR x=10) x++; • Solution: Try: if(x>0 AND x<=10) x++; 4 19 Logic and e-Mail Filtering I • Situation: You are tired of getting spam about losing weight and making money on eBay.
• Problem: You tried the following filter: If(subject contains weight and subject contains eBay) Delete message • Solution: You meant: If(subject contains weight or subject contains eBay) Delete message 20 Logic in Programming II • Situation: Consider the following loop: while(NOT(A[i]!= 0 AND NOT(A[i]>= 10))) • Problem: You are convinced this is way too complicated.
• Solution: Well, it’s hard to say at this point… 21 Logic in Programming III • Situation: Consider the following loop: while( (i10) OR (i<size AND A[i]<0) OR NOT (A[i]!= 0 AND NOT (A[i]>= 10) ) ) • Problem: You are convinced this is way too complicated.
• Solution: Yet another example we can’t solve.
22 Logic and e-mail Filtering II • Situation: You get so much spam, you decide to delete any message not sent (or copied) to you (blah), or a group to which you belong (foo). You try: if( TO or CC does not contain blah OR TO or CC does not contain foo) Delete message • Problem: It seems like all of your e-mail is being deleted.
• Solution: This one is confusing, and we will solve it later. 23 Logic and Medication • Situation: Your new medication has the following label: Take 1 or 2 pills every 4-6 hours until condition improves. Do not exceed 6 pills per day, or take for longer than 7 days unless directed by a doctor. Do not take this drug during the last 3 months of pregnancy, unless directed by a doctor. If you are taking XYZs and PDQs or have taken either an XYZ or a PDQ within the last 90 days, and the other within the last 30 days, or are taking an ABC or have taken an ABC within the last 60 days, you should not take this drug. Do not drink alcohol or smoke while on this drug. If you have a heart condition, asthma, diabetes, or have an IQ below 25, do not take this drug. Do not take this drug if you have a high fever, cold sweats, runny nose, headache, or sore throat, and discontinue use if you develop any of these symptoms, and dizziness, nervousness, or sleeplessness occur.
• Problem: Under what conditions can you take it?
• Solution: Well, this one may take more than simple logic.
24 Why Logic?
• Hopefully the last several examples have convinced you that knowing more logic is important for computer science, and life in general.
• If you are still not convinced, tough. You are going to learn it anyway.
5 25 Propositional Equivalences • Many problems, included the last several examples, can be solved by understanding the concept of propositional equivalences.
• Example: The statement “I am not a student and I am not living in the dorm” is equivalent to “I am not a student or living in the dorm.” • Example: “You will pass this class or you will not be in the J.D. Edwards program next year” is equivalent to “If you are to be in the J.D. Edwards program next year, then you must pass this class.” 26 Some Terminology • Definition: A tautology is a proposition that is always true.
• Definition: A contradiction is a proposition that is always false.
• Definition: A proposition that is not a tautology or a contradiction is a contingency.
27 Propositional Equivalence • Definition #1: Propositions p and q are called logically equivalent if p↔q is a tautology. • Definition #2: Propositions p and q are logically equivalent if and only if they have the same truth table.
• Notation: If p and q are equivalent, we write p⇔q • Example: The propositions ¬p∨q and p→q are logically equivalent. We can see this by constructing the truth tables.
T F F F F T T T q p T T F T ¬p∨q T T F F ¬p T T F T p→q 28 Another Example • We show that p∨(q∧r) and (p∨q)∧(p∨r) are logically equivalent.
F F F F T T T T p F F T F F T T T F F T F F T T T r q F F F T T T T T (p∨q)∧(p∨r) F T F T T T T T p∨r F F T T T T T T p∨q F F F T T T T T p∨(q∧r) F F F T F F F T q∧r 29 Simple Logical Equivalences Equivalences involving one proposition Name Equivalence p∧F ⇔F Domination laws p∨T ⇔T p∨F ⇔p Identity laws p∧T ⇔p p∧p ⇔p Idempotent laws p∨p ⇔p Double negation law ¬(¬p) ⇔p (Not an offical name) p∧¬p ⇔F Cancellation laws p∨¬p ⇔T 30 Logical Equivalences Equivalences involving multiple propositions Name Equivalence p∧q ⇔q∧p Commutative laws p∨q ⇔q∨p (p∧q)∧r ⇔p∧(q∧r) Associative laws (p∨q)∨r ⇔p∨(q∨r) p∧(q∨r) ⇔(p∧q)∨(p∧r) Distributive laws p∨(q∧r) ⇔(p∨q)∧(p∨r) ¬(p∨q) ⇔¬p∧¬q De Morgan’s laws ¬(p∧q) ⇔¬p∨¬q Implication law (p→q) ⇔(¬p∨q) 6 31 Using Logical Equivalences I Example 1: • Show that (p∧q)→q is a tautology using logical equivalences.
q q p ∨ ¬ ∨ ¬ ⇔ ) ( ←Implication law ←De Morgan’s Law ←Associative law ←Cancellation Law ←Domination Law ) ( ) ) (( q q p q q p ∨ ∧ ¬ ⇔ → ∧ ) ( q q p ∨ ¬ ∨ ¬ ⇔ T ∨ ¬ ⇔ p T ⇔ 32 Using Logical Equivalences II Example 2: Show that (¬(q→p))∨(p∧q) is logically equivalent to q ↓Implication law ←De Morgan’s and double negation ←Commutative law ←Distributive law ) ( )) ( ( ) ( )) ( ( q p p q q p p q ∧ ∨ ∨ ¬ ¬ ⇔ ∧ ∨ → ¬ ) ( ) ( q p p q ∧ ∨ ¬ ∧ ⇔ ) ( ) ( p q p q ∧ ∨ ¬ ∧ ⇔ ) ( p p q ∨ ¬ ∧ ⇔ T ∧ ⇔q q ⇔ ←Cancellation law ←Identity law 33 Logic in Programming II • Situation: Consider the following loop: while(NOT(A[i]!=0 AND NOT(A[i]>= 10))) • Problem: You are still convinced this is way too complicated, and now you think you can simply it.
• Solution: We can use De Morgan’s law and the double negation law to obtain while( A[i]==0 OR A[i]>= 10) 34 Logic in Programming III • Situation: Consider the following loop: while( (i10) OR (i<size AND A[i]<0) OR NOT (A[i]!= 0 AND NOT (A[i]>= 10))) • Problem: You are convinced this is way too complicated, and with some work, you can simplify it.
• Solution: Start by simplifying the last part is in the last example: while( (i10) OR (i<size AND A[i]<0) OR (A[i]==0 OR A[i]>= 10) ) • Then, use the distributive law: while( (i10 OR A[i]<0) ) OR (A[i]==0 OR A[i]>= 10) ) 35 An Important Note • In many programming languages, including Java, C++, and C, applying the commutative law to a proposition may or may not be a good idea.
• The reason for this is that these languages use a technique sometimes call “short circuiting.” • For instance, if A is an array of n elements, the statements if(i<n AND A[i]==0) and if(A[i]==0 AND i<n) are NOT equivalent. Why?
36 Logic and e-Mail Filtering II • Situation: You get so much spam, you decide to delete any message not sent (or copied) to you (blah), or a group to which you belong (foo). You try: if( TO or CC does not contain blah OR TO or CC does not contain foo) Delete message • Problem: It seems like all of your e-mail is being deleted.
• Solution: This one is a little more complicated. We start by applying De Morgan’s law: if(NOT (TO or CC contains blah AND TO or CC contains foo) ) Delete message • Let p=“TO or CC contains blah” and q=“TO or CC contains foo”. 7 37 Logic and e-Mail Filtering II • The statement becomes if(NOT (p AND q)) Delete message • The truth table for NOT(p AND q)is: F T F F F F F F T T T T p∧q q p T T T F ¬(p∧q) • So e-mail is deleted unless TO or CC contain both blah and foo, which is clearly not what we wanted.
• If TO or CC contains either blah or foo, we do not want to delete. The truth table we want is… T F F F ?
F T T T p∨q • Applying negation, we now recognize this as: 38 Logic and e-Mail Filtering II • Then what we really want is: if(NOT(p OR q)) Delete message F T F T F F F F T F T T ¬(p∨q) q p •Using DeMorgan’s Law, it becomes: if(NOT p AND NOT q)) Delete message •Retranslating, we seem to have wanted if( TO or CC does not contain blah AND TO or CC does not contain foo) Delete message •Great. We did it.
•Or did we?
39 Logic and e-Mail Filtering II • We now think that the following filter should work: if( TO or CC does not contain blah AND TO or CC does not contain foo) Delete message • Unfortunately, all of your e-mail is still being deleted.
• Let’s keep trying. Let p = “TO contains blah” q = “CC contains blah” r = “TO contains foo” s = “CC contains foo” • The filter is if( (NOT p OR NOT q ) AND (NOT r OR NOT s) ) Delete message • This is where we went wrong earlier… 40 Logic and e-Mail Filtering II • Continuing, we had if( (NOT p OR NOT q ) AND (NOT r OR NOT s) ) Delete message • Applying De Morgan’s Law, we have if( NOT (p AND q ) AND NOT( r AND s) ) Delete message • Applying De Morgan’s Law again, we have if( NOT( (p AND q ) OR ( r AND s) ) ) Delete message • In English, the condition is loosely translated “If it is not the case that either both TO AND CC contains blah OR both TO AND CC contains foo” 41 Logic and e-Mail Filtering II • In other words, unless the person sending the message put your address in both the TO and CC fields, it will be deleted.
• The problem is that the filter parses “if TO OR CC (X)” as “if TO (X) OR CC (X).” • This means that “if TO OR CC (NOT X)” is parsed as “if TO (NOT X) OR CC (NOT X),” not as “if NOT ( TO (X) OR CC (X)),” which is what we did earlier in the example. • Because of this, you should not use the “if TO OR CC…” filter with the “does not contain” condition, since it is most likely not what you intended.
42 Logic and e-Mail Filtering II • Now that we know the problem, we can fix it.
• What we want is really: If(TO does not contain blah AND CC does not contain blah AND TO does not contain foo AND CC does not contain foo) Delete message • Recall the original filter was: if(TO or CC does not contain blah OR TO or CC does not contain foo) Delete message 8 43 Logic and e-Mail Filtering II • This example should help illustrate the following 1. Sometimes we can think a statement means one thing when it actually means another.
2. Sometimes, we simply can’t figure out what a statement means at all.
• In these cases, we can use logic to assist us in determining the true meaning of statements.
44 Logic and Medication • Situation: Your new medication has the following label: Take 1 or 2 pills every 4-6 hours until condition improves. Do not exceed 6 pills per day, or take for longer than 7 days unless directed by a doctor.
Do not take this drug during the last 3 months of pregnancy, unless directed by a doctor. If you are taking XYZs and PDQs or have taken either an XYZ or a PDQ within the last 90 days, and the other within the last 30 days, or are taking an ABC or have taken an ABC within the last 60 days, you should not take this drug. Do not drink alcohol or smoke while on this drug. If you have a heart condition, asthma, diabetes, or have an IQ below 25, do not take this drug. Do not take this drug if you have a high fever, cold sweats, runny nose, headache, or sore throat, and discontinue use if you develop any of these symptoms, and dizziness, nervousness, or sleeplessness occur.
• Problem: Under what conditions can you take it?
• Solution: O.K., I give up. But seriously, there may be a time when you really need to solve a similar problem. 45 Some Exercises 1. Construct the truth table for the following propositions a) (p→q)∧p b) (¬p↔¬q) ↔(q↔r) c) p∧(q→¬ (r ∨q)) 2. I do not want any e-mail that contains the words puke, ralph, or hurl, unless it was specifically sent to me (blah). How do I do it?
3. Is ¬(p→q)→¬q a tautology? Give two different proofs.
4. Show that p↔q and (p∧q)∨(¬p∧¬q) are logically equivalent.
5. Show that [(p∨q)∧(p→r)∧(q→r)]→r is a tautology. Give a proof using equivalences and a truth table.
46 6.
Can you enter this contest? Can I? Can your friends?
Rules: Must be at least 18 years old to enter. Must be enrolled in a computer science course at NU, but cannot be enrolled in a course in the J.D. Edwards Program. Full-time employees of UNL, UNO, and UNK are not allowed to enter, unless they are faculty of the computer science department, or work in information systems support. 7.
Your mom says “If you don’t eat your meat, you can’t have any pudding.” What should you do? Why?
8.
Your friend offers “Heads I win, tails you lose.” Do you take the bet?
9.
If you can’t understand this problem, then you didn’t learn the material from these notes. Did you learn the material from these notes?
10.
If you don’t want to not learn the material from this class, you should not fail to not skip doing problems. Or should you?
11.
In an episode of The Simpsons Bart said something similar to “All I know is that I’m getting straight A’s, and that ain’t not bad.” Does he deserve all A’s? |
13733 | https://www.youtube.com/watch?v=sbAemzOzlww | 20% of a Number Equals 3 Times the Number Plus 28 — What’s the Number?
TabletClass Math
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In this video, I’ll solve an interesting percent and algebra problem:
20% of a number equals three times that number plus 28. What’s the number?
This is a great example of how algebra and percent skills work together in real problems — and a perfect warm-up for GED, SAT, or math review.
If this helps you out, be sure to like, subscribe, and check out the links below for full lessons and practice!
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a lot of people are going to be confused on how to solve this math problem let's take a look at it so 20% of a number is 3 the number plus 28 what is this number well I'm going to show you exactly how to solve this but before we get started my name is John and I have been teaching math for decades and if you need help in math make sure to check out my math program at tcmathacademy.com you can find the link in the description all right so let's go ahead and get into this right now okay so first things first first we are dealing with a math word problem so always uh use the rule of three which is even though you understand the problem uh always force yourself to read the problem at least three times and just don't read it fast like okay here's the problem here's the problem here's the problem that's not the point that's not the rule of three that's not what I'm saying in other words read the prom one time just to get your bearings read the prom slow down read it again and as you're reading it you should be thinking about all right have I seen a problem like this before uh well I have to use algebra to solve this problem uh you know you know you're thinking about some of the things that you have in your math toolkit that can help you solve this problem and then the third time you read that problem and this is the minimum amount of times you really want to make sure you understand the question and once you understand the question you want to circle back and kind of connect the question to the information in the problem so here we are dealing with percent so we're going to have to know something about percent and we're talking about some number a number well what is a number well we don't know what a number is so uh using a variable like x is a good uh you know idea because a variable um of course when we're talking about variables we're talking about algebra a variable can represent a number so we again we want to translate this verbal situation into a mathematical situation so a lot of different skills going on here not difficult skills but again things that hopefully you learn uh you know in your basic uh pre-alggebra or algebra 1 courses all right so now that we have uh a good idea what this problem is about we want to translate or model the problem and again we're going to want to use a variable because we're dealing with an unknown value see when you solve a math problem and there's some sort of unknown value that's in the problem you oftentimes want to think of using algebra right now some of you may have been creative enough to kind of solve this problem on your own without algebra and that you know is uh to be applauded but algebra is a tool and it just makes things so much easier all right so I'm going to go ahead and define this number or mystery number i'm going to let the variable x equal that number right the number so uh if I solve for x well then I will have solved the problem okay so you want to kind of delineate these things as well in other words just don't say x or whatever the case is you should define it hey this x is going to equal this mystery number then we can go back to our problem and start building uh you know some relationship between the information in the problem now namely what you want to do is try to build an equation because you can't solve for a variable like x unless we have some sort of equation and we want to use the information in the problem to build an equation so let's go ahead and do that right now and again uh this uh part of the problem is going to depend upon your ability to translate a verbal phrase i.e a phrase with uh letters okay uh sentences into um algebra okay into a variable phrase so let's go ahead and break this into two separate components so the first part is 20% of a number now I have this broken up into two separate parts and you'll see why here in a second but anytime you see the word is that is a a very important word uh in mathematics the is word the the word is is the equal sign okay so we just you know instantaneously be like oh there's is so you want to put that equal sign so we're equating something you know we're equating before everything before the word is to everything that's after the word is okay so 20% of a number is or is equal to 3 the number plus 28 so now we want to go ahead and translate each one of these things all right so 20% of a number now remember we defined that number or we're going to let our number equal to x but we need to express this uh in terms of uh you know mathematics so 20% of a number well what does it mean to find the percent of a number well if you forget uh forget uh go back and just do something easy use an easy example to help you remember how to do this stuff so let's just take something like 20% of 100 so what's 20% of 100 don't even look at the math what is 50% of 100 okay hopefully you're saying "Hey Mr youtube Math Man that's 50 what's 100% of 100?" Well it would be 100 what's 1% of 100 it would be one so 20% of 100 is obviously 20 okay so you know that right so let's say you forgot how to do this you're like I think this is how you find the percent of a number so just do take a problem that you already know the answer to and then uh review what you think is the procedure now the procedure to find the percent of a number is to uh convert that percent to a decimal now how do we do that well we divide that by 100 so we're going to divide our percent by 100 so 20% uh divided by 100 is equal to 0.2 now another way you can think of that is you can uh move the decimal point over two places to the left uh that is just a result of dividing by 100 so some of you might be saying hey Mr youtube 2 math man just move the decimal point over two places to the left yes indeed uh you know I plan on doing that but why are we doing that well that's just the result of dividing by 100 okay so that's what we're going to do we're going to take our percent and turn it into a decimal and then just multiply by the number that we want to find that percent of all right so 20% is 02 so 0.2 100 again I said feel free to use your calculator but 0.2 100 is indeed 20 all right so how do we find the percent of a number well we're going to convert or write that percent as a decimal and multiply by that number so in this case this would be 2 x all right so that's going to be the um translation for this part of the problem now the second part of the problem is what well this is 3 the number well that's going to be 3 what x plus 28 all right so hopefully this is all making sense and now let's put it all together so 20% of our number is going to be 2 uh 0.2 times that number which of course we don't know it's x right so this is 20% of the number is equal to 3 that number plus 28 all right so this is what we have here we have this lovely a linear equation 2x = 3x + 28 now we can put our attention on solving this equation for x all right so that is what we're going to do now real quick if you want my best math instruction you definitely got to check out my full courses again you can find links to these in the description of this video but they span basic math to advanced math and everything in between okay so let's keep going with this problem and don't forget to like and subscribe as that definitely helps me out okay so now we're going to solve for x we have 2x equal to 3x + 28 and here is all the steps we're going to go ahead and subtract 3x from both sides of the equation we're going to move all our variables to the left and when we do this we're going to be very careful of course right 2 minus uh uh 3 is negative you got to be careful here right because we're dealing with positive and negative numbers -2.8x is equal to 28 all right so we have 2.8x is equal to 28 so now to solve for x all we have to do is divide both sides of the equation by -2.8 and when we do that we get 28 / -2.8 x is going to be equal to -10 all right now let's suppose we did all this and we're saying you know what we're feeling pretty good about this answer but maybe you might be saying well maybe I'm just not sure okay how could we check this answer well this is not that difficult matter of fact let's go back and say all right well our number right uh what we were looking for x is the number is -10 all right so this number is -10 let's go back and check this uh translation so 20% of our number now we know the number is -10 is what well 20% of -10 is going to be2 -10.2 -10 is -2 all right so let's see if that is equal to this part of the problem 3 the number + 28 well what is the number well remember the number now is -10 so 3 the number would be 3 -10 3 -10 is -30 + 28 -30 + 28 is indeed -2 wow look at this everything worked out okay so a simple example but an important uh example on you know basic skills that you need in order to solve an algebra word problem and namely you need to obviously know how to solve equations but you need to know how to translate verbal phrases into algebraic or variable phrases and if anything in this uh you know prom if you don't understand just make yourself a math shopping list be like oh I need to pick up some skills on this this this and this this is what I'm talking talking about you can improve in math but you can't improve if you're not willing to take action now some of you out there are um you know not math students but you want to relearn math check out my two courses you can find links to these in the description first is my math foundation course that's just a quick basic math review decimals fractions all that really super important stuff that you have to get down but if you're um kind of wanting to relearn math let's say you've been away from math for a long time well then check out my math skills rebuilder course in that course I teach you basic math algebra geometry some basic trigonometry and some probability and statistics and all my courses are self-paced as well okay so with all that being said I definitely wish you all the best in your math adventures thank you for your time and have a great day |
13734 | https://www.khanacademy.org/science/electrical-engineering/ee-circuit-analysis-topic/ee-resistor-circuits/a/ee-delta-wye-resistor-networks | Delta-Wye resistor networks (article) | Khan Academy
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Electrical engineering
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Lesson 2: Resistor circuits
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Delta-Wye resistor networks
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The Delta-Wye transformation is an extra technique for transforming certain resistor combinations that cannot be handled by the series and parallel equations. This is also referred to as a Pi - T transformation.
Sometimes when you are simplifying a resistor network, you get stuck. Some resistor networks cannot be simplified using the usual series and parallel combinations. This situation can often be handled by trying the Δ−Y_transformation_, or 'Delta-Wye' transformation.
The names _Delta_ and _Wye_ come from the shape of the schematics, which resemble letters. The transformation allows you to replace three resistors in a Δ configuration by three resistors in a Y configuration, and the other way around.
The Δ−Y drawing style emphasizes these are 3-terminal configurations. Something to notice is the different number of nodes in the two configurations. Δ has three nodes, while Y has four nodes (one extra in the center).
The configurations can be redrawn to square up the resistors. This is called a π−T configuration,
The π−T style is a more conventional drawing you would find in a typical schematic. The transformation equations developed next apply to π−T as well.
Δ−Y transformation
For the transformation to be equivalent, the resistance between each pair of terminals must be the same before and after. It is possible to write three simultaneous equations to capture this constraint.
Consider terminals x and y (and for the moment assume terminal z isn't connected to anything, so the current in R 3 is 0 ). In the Δ configuration, the resistance between x and y is R c in parallel with R a+R b .
On the Y side, the resistance between x and y is the series combination R 1+R 2 (again, assume terminal z isn't connected to anything, so R 1 and R 2 carry the same current and can be considered in series). We set these equal to each other to get the first of three simultaneous equations,
R 1+R 2=R c(R a+R b)R c+(R a+R b)
We can write two similar expressions for the other two pairs of terminals. Notice the Δ resistors have letter names, (R a, etc.) and the Y resistors have number names, (R 1, etc.).
After solving the simultaneous equations (not shown), we get the equations to transform either network into the other.
Δ→Y transformation
Equations for transforming a Δ network into a Y network:
R 1=R b R c R a+R b+R c
R 2=R a R c R a+R b+R c
R 3=R a R b R a+R b+R c
Transforming from Δ to Y introduces one additional node.
Y→Δ transformation
Equations for transforming a Y network into a Δ network:
R a=R 1 R 2+R 2 R 3+R 3 R 1 R 1
R b=R 1 R 2+R 2 R 3+R 3 R 1 R 2
R c=R 1 R 2+R 2 R 3+R 3 R 1 R 3
Transforming from Y to Δ removes one node.
Example
Let's do a symmetric example. Assume we have a Δ circuit with 3 Ω resistors. Derive the Y equivalent by using the Δ→Y equations.
R 1=R b R c R a+R b+R c=3⋅3 3+3+3=1 Ω
R 2=R a R c R a+R b+R c=3⋅3 3+3+3=1 Ω
R 3=R a R b R a+R b+R c=3⋅3 3+3+3=1 Ω
Going in the other direction, from Y→Δ, looks like this,
R a=R 1 R 2+R 2 R 3+R 3 R 1 R 1=1⋅1+1⋅1+1⋅1 1=3 Ω
R b=R 1 R 2+R 2 R 3+R 3 R 1 R 2=1⋅1+1⋅1+1⋅1 1=3 Ω
R c=R 1 R 2+R 2 R 3+R 3 R 1 R 3=1⋅1+1⋅1+1⋅1 1=3 Ω
Example
Now for an example that's a little less tidy. We want to find the equivalent resistance between the top and bottom terminals.
Try as we might, there are no resistors in series or in parallel. But we are not stuck. First, let's redraw the schematic to emphasize we have two Δ connections stacked one on the other.
Now select one of the Δ's to convert to a Y. We will perform a Δ→Y transformation and see if it breaks the logjam, opening up other opportunities for simplification.
We go to work on the bottom Δ (an arbitrary choice). _Very carefully_ label the resistors and nodes. To get the right answers from the transformation equations, it is critical to keep the resistor names and node names straight. R c must connect between nodes x and y, and so on for the other resistors. Refer to _Diagram_ _1_ above for the labeling convention.
When we perform the transform on the lower Δ, the black Δ resistors will be replaced by the new gray Y resistors, like this:
Perform the transform yourself before looking at the answer. Check that you select the right set of equations.
Compute three new resistor values to convert the Δ to a Y, and draw the complete circuit.
Show the Delta to Y transformation.
Apply the transformation equations for Δ→Y.
R 1=R b R c R a+R b+R c=5⋅3 4+5+3=15 12=1.25 Ω
R 2=R a R c R a+R b+R c=4⋅3 4+5+3=12 12=1 Ω
R 3=R a R b R a+R b+R c=4⋅5 4+5+3=20 12=1.66 Ω
Derive an equivalent Y network by substituting the Δ resistors. Make sure the Y resistor names connect between the proper node names. Refer to _Diagram_ _1_ above for the labeling convention.
And voilà! Check out our circuit. It now has series and parallel resistors where it had none before. Continue simplification with series and parallel combinations until we get down to a single resistor between the terminals. Redraw the schematic again to square up the symbols into a familiar style.
We proceed through the remaining simplification steps just as we did before in the article on Resistor Network Simplification.
On the left branch, 3.125+1.25=4.375 Ω
On the right branch, 4+1=5 Ω
The two parallel resistors combine as 4.375||5=4.375⋅5 4.375+5=2.33 Ω
And we finish by adding the last two series resistors together,
R e q u i v a l e n t=2.33+1.66=4 Ω
Summary
Δ−Y transformations are another tool in your bag of tricks for simplifying circuits prior to detailed analysis.
Don't memorize the transformation equations. If the need arises, you can look them up.
Copyright notice
This article is licensed under CC BY-NC-SA 4.0.
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Fasteric Algorithm 9 years ago Posted 9 years ago. Direct link to Fasteric Algorithm's post “Why we can assume that te...” more Why we can assume that terminal z isn't connect to anything?
I mean, in the example, it clearly show that terminal z did connect to something and must have a current flow through.
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APDahlen 9 years ago Posted 9 years ago. Direct link to APDahlen's post “Hello Fasteric, Try to t...” more Hello Fasteric,
Try to think of each circuit as a black box - assume you don't know what is inside. When you "look" from terminal A to B you see a certain resistance. When you "look" from terminal B to C you see a certain resistance. Finally when you "look" from terminal C to A you again see a resistance. It's important that you "look" at the circuit in isolation. If it was connected to other things it would be difficult to do this thought experiment.
Please give it a try you should be able to convince yourself that the wye and delta look the same. If it helps try calculating the resistances when the delta has three resistors each with a value of 3 Ohm. Then try the wye with three resistors each having a 1 Ohm value.
Regards,
APD
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popalzaymaiwand 2 years ago Posted 2 years ago. Direct link to popalzaymaiwand's post “A video explanation might...” more A video explanation might made it easier to understand.
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Hagar Usama 8 years ago Posted 8 years ago. Direct link to Hagar Usama's post “I can derive R1, R2 & R3 ...” more I can derive R1, R2 & R3 easily, but I can't derive Ra at all. I can't figure it out !
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heatdeath 2 years ago Posted 2 years ago. Direct link to heatdeath's post “I found this lesson to be...” more I found this lesson to be harder to follow that usual.
Here are some ways to improve it:
1 Put the circuit from the example right at the top. I lost concentration my first time through the lesson because there were a lot of equations and it wasn't clear what it was all for.
2 Don't skip over the general formulas for the parallel and series resistances. We haven't had to use them much in the lesson so far and we still aren't familiar with them.
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Daniel 8 years ago Posted 8 years ago. Direct link to Daniel's post “1) Why in the end we get ...” more 1) Why in the end we get 4 Ohms?
If we assume, that Z isn't connected to any source, then we can put resistor = 1.66 Ohms in the circuit. In this case resistor = 4.375 and resistor = 1.66 Ohms are going to be in series. Sum them up = 6.035 Ohms.
We're left with resistor = 6.035 Ohms on the left and resistor = 5 Ohms on the right. They are in parallel -> apply the formula -> (6.0355)/(6.035+5) = 2,73.
2) Paraller resistors theory says, that resistors are in parallel, if they both share same node and same voltage. Assume we've connected voltage source to the top and we have 2V. So the voltage which goes into resistor = 5 Ohms is different from voltage which goes into resistor = 4.375 Ohms ( cuz voltage also goes into resistor = 1.66 Ohms ). And they also don't share the same note, since there is resistor = 1.66 of the way. How does it work?
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Mirk a year ago Posted a year ago. Direct link to Mirk's post “Are there a whole lot of ...” more Are there a whole lot of transformation equations, or are there several basic ones that can be filled in by solving many equations?
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Willy McAllister a year ago Posted a year ago. Direct link to Willy McAllister's post “For resistor circuits the...” more For resistor circuits the basic equations are Ohm's Law, v=ir, series resistors Rs=R1+R2+..., and parallel resistors, 1/Rp=1/R1+1/R2... These three solve 99% of resistor simplification problems. The last 1% are the rare cases where series/parallel does not work, and are solved by the Delta-Wye transform equations described here.
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Juniper 2 years ago Posted 2 years ago. Direct link to Juniper's post “Why is Rc in parallel wit...” more Why is Rc in parallel with Ra+Rb? It has not been explained where this comes from. (After watching the next video, it seems that the voltage divider video would have been helpful before this topic, using the two assumptions of zero current through the x,y, or z to view the circuit in different ways.
Am I thinking about this right?)
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Willy McAllister 2 years ago Posted 2 years ago. Direct link to Willy McAllister's post “You are thinking about th...” more You are thinking about this correctly. When you analyze the resistance between any two nodes you assume the third node is open. (In the section titled D - Y Transformation the assumption about node z being open is mentioned in two places.)
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Francis Billones 5 years ago Posted 5 years ago. Direct link to Francis Billones's post “Hey, I don't understand w...” more Hey, I don't understand why there aren't any series or parallel resistors in the first circuit. I can find 2 parallel resistors there.
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omeva.2010 5 years ago Posted 5 years ago. Direct link to omeva.2010's post “Where can I learn more in...” more Where can I learn more in detail about THREE PHASE CIRCUITS ?
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Wild Chaser 9 years ago Posted 9 years ago. Direct link to Wild Chaser's post “how are r1 and r2 in seri...” more how are r1 and r2 in series in the y-transformation in the 1st diagram? since there is a branch of resistor in the common node
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13735 | https://schools.aglasem.com/ncert-books-class-10-maths-chapter-5/ | NCERT Book Class 10 Maths Chapter 5 Arithmetic Progressions (PDF) - AglaSem
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Home » 10th Class » NCERT Book Class 10 Maths Chapter 5 Arithmetic Progressions (PDF)
NCERT Book Class 10 Maths Chapter 5 Arithmetic Progressions (PDF)
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NCERT Book Class 10 Maths Chapter 5 Arithmetic Progressions is here. You can read and download Class 10 Maths Chapter 5 PDF from this page of aglasem.com. Arithmetic Progressions is one of the many lessons in NCERT Book Class 10 Maths in the new, updated version of 2023-24. So if you are in 10th standard, and studying Maths textbook (named Mathematics), then you can read Ch 5 here and afterwards use NCERT Solutions to solve questions answers of Arithmetic Progressions.
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NCERT Book Class 10 Maths Chapter 5 Arithmetic Progressions
The complete Chapter 5, which is Arithmetic Progressions, from NCERT Books for Class 10 Maths is as follows.
NCERT Book Class 10 Maths Chapter 5 Arithmetic Progressions PDF Download Link – Click Here To Download The Complete Chapter PDF
NCERT Book Class 10 Maths Full Book PDF Download Link – Click Here To Download The Complete Book PDF
NCERT Book Class 10 Maths Chapter 5 Arithmetic Progressions PDF
The direct link to download class 10 Maths NCERT Book PDF for chapter 5 Arithmetic Progressions is given above. However if you want to read the complete lesson on Arithmetic Progressions then that is also possible here at aglasem. So here is the complete class 10 Maths Ch 5 Arithmetic Progressions.
NCERT Book Class 10 Maths Chapter 5 Arithmetic Progressions ViewDownload
NCERT Book for Class 10 Maths
Besides the chapter on Arithmetic Progressions, you can read or download the NCERT Class 10 Maths PDF full book from aglasem. Here is the complete book:
Chapter 1: Real Numbers
Chapter 2: Polynomials
Chapter 3: Pair Of Linear Equations In Two Variables
Chapter 4: Quadratic Equations
Chapter 5: Arithmetic Progressions
Chapter 6: Triangles
Chapter 7: Coordinate Geometry
Chapter 8: Introduction to Trigonometry
Chapter 9: Some Applications of Trigonometry
Chapter 10: Circles
Chapter 11: Area related to circles
Chapter 12: Surface Areas And Volumes
Chapter 13: Statistics
Chapter 14: Probability
NCERT Books for Class 10
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NCERT Book Class 10 English
NCERT Book Class 10 Hindi
NCERT Book Class 10 Maths
NCERT Book Class 10 Sanskrit
NCERT Book Class 10 Science
NCERT Book Class 10 Social Science
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Class 10 Maths Chapter 5 Arithmetic Progressions NCERT Textbook – An Overview
The highlights of this Arithmetic Progressions chapter PDF are as follows.
| Aspects | Details |
--- |
| Class | 10 |
| Subject | Maths |
| Book | Mathematics |
| Chapter Number | Ch 5 |
| Chapter Name | Arithmetic Progressions |
| Book Portion Here | NCERT Book Class 10 Maths Ch 5 Arithmetic Progressions |
| Download Format | PDF |
| Version | NCERT Book (New, Updated) 2023-24 |
| Complete Book | NCERT Book Class 10 Maths |
| All Class 10 Books | NCERT Books for Class 10 |
| All Textbooks | NCERT Books |
| NCERT Books in Hindi | NCERT Books for Class 10 in Hindi |
| NCERT Solutions | NCERT Solutions for Class 10 |
| More Study Material | Class 10 |
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13736 | https://www.uhhospitals.org/blog/articles/2014/04/best-and-worst-foods-for-acid-reflux | The Best and Worst Foods for Acid Reflux – What to Eat and Avoid | University Hospitals
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The Best and Worst Foods for Acid Reflux
April 15, 2014
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A hot burning in the chest, a bitter taste in the throat, a gassy bloating in the stomach – acid reflux is no picnic. What you eat, however, can have an impact. The best and worst foods for acid reflux could spell the difference between sweet relief and sour misery.
What Aggravates Acid Reflux?
Acid reflux occurs when the sphincter at the base of the esophagus isn’t working well, allowing fluid from the stomach to enter the esophagus. The worst foods for reflux can worsen painful symptoms, while other foods can soothe them, says UH gastrointestinal surgeon Leena Khaitan, MD
“Diet changes can significantly affect acid reflux and allow you to avoid other treatments,” Dr. Khaitan says.
Best Foods for Acid Reflux
“A diet balanced with vegetables, protein and fruits is best,”Dr. Khaitan says. Examples of the best foods for acid reflux include:
Chicken breast– Be sure to remove the fatty skin. Skip fried and instead choose baked, broiled or grilled.
Lettuce, celery and sweet peppers –These mild green veggies are easy on the stomach – and won’t cause painful gas.
Brown rice– This complex carbohydrate is mild and filling – just don’t serve it fried.
Melons– Watermelon, cantaloupe and honeydew are all low-acid fruits that are among the best foods for acid reflux.
Oatmeal– Filling, hearty and healthy, this comforting breakfast standard also works for lunch.
Fennel– This low-acid crunchy vegetable has a mild licorice flavor and a natural soothing effect.
Ginger– Steep caffeine-free ginger tea or chew on low-sugar dried ginger for a natural tummy tamer.
Worst Foods for Reflux
In general, anything that is fatty, acidic or highly caffeinated should be avoided. The worst foods for acid reflux list includes:
Coffee and tea– Caffeinated beverages aggravate acid reflux. Opt for teas without caffeine.
Carbonated beverages– The bubbles expand in your stomach, creating more pressure and pain. Choose plain water or decaf iced tea.
Chocolate– This treat has a trifecta of acid reflux problems: caffeine, fat and cocoa.
Peppermint–Don’t be fooled by its reputation for soothing the tummy; peppermint is an acid reflux trigger.
Grapefruit and orange– The high acidity of citrus fruits relaxes the esophagus sphincter and worsens symptoms.
Tomatoes– Also avoid marinara sauce, ketchup and tomato soup – they’re all naturally high in acid.
Alcohol–This has a double whammy effect. Alcohol relaxes the sphincter valve but it also stimulates acid production in the stomach.
Fried foods –These are some of the worst foods for reflux. Skip the french fries, onion rings and fried chicken — cook on the grill or in the oven at home.
Late-night snacks– Avoid eating anything in the two hours before you go to bed. Also, you can try eating four to five smaller meals throughout the day instead of two to three large meals.
When to Talk to Your Doctor About Acid Reflux
It's a good idea to speak with your doctor if the best foods for acid reflux do not relieve your symptoms, Dr. Khaitan says. Other options can include lifestyle changes, medications to block acid, and surgical procedures on the esophagus sphincter.
It is important to make a doctor's appointmentif you have heartburn or acid reflux that is severe or frequent, Dr. Khaitan adds. Chronic acid reflux is known as gastroesophageal reflux disease (GERD) and can lead to esophageal cancer.
Related Links
University Hospitals’ experienced team of digestive health specialistsoffers innovative diagnostics and treatments for all stages of GERD.
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Tags: GERD, Esophageal Disease, Men’s Health, Men’s Health: Full-Life Fitness
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Does frequent heartburn come between you and your favorite foods at outdoor cookouts? It doesn’t have to - here are some common sense tips to help...
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13737 | https://activecalculus.org/multi/S-11-4-Double-Integrals-Applications.html | Applications of Double Integrals
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Active Calculus - Multivariable
Steve Schlicker, Mitchel T. Keller, Nicholas Long
ContentsIndex
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Front Matter
Colophon
Features of the Text
Vector Calculus Preface
Acknowledgments
Active Calculus - Multivariable: our goals
How to Use this Text
9 Multivariable and Vector Functions
9.1 Functions of Several Variables and Three Dimensional Space
9.1.1 Functions of Several Variables
9.1.2 Representing Functions of Two Variables
9.1.3 Some Standard Equations in Three-Space
9.1.4 Traces
9.1.5 Contour Maps and Level Curves
9.1.6 A gallery of functions
9.1.7 Summary
9.1.8 Exercises
9.2 Vectors
9.2.1 Representations of Vectors
9.2.2 Equality of Vectors
9.2.3 Operations on Vectors
9.2.4 Properties of Vector Operations
9.2.5 Geometric Interpretation of Vector Operations
9.2.6 The Magnitude of a Vector
9.2.7 Summary
9.2.8 Exercises
9.3 The Dot Product
9.3.1 The Dot Product
9.3.2 The angle between vectors
9.3.3 The Dot Product and Orthogonality
9.3.4 Work, Force, and Displacement
9.3.5 Projections
9.3.6 Summary
9.3.7 Exercises
9.4 The Cross Product
9.4.1 Computing the cross product
9.4.2 The Length of u×v
9.4.3 The Direction of u×v
9.4.4 Torque is measured by a cross product
9.4.5 Comparing the dot and cross products
9.4.6 Summary
9.4.7 Exercises
9.5 Lines and Planes in Space
9.5.1 Lines in Space
9.5.2 The Parametric Equations of a Line
9.5.3 Planes in Space
9.5.4 Summary
9.5.5 Exercises
9.6 Vector-Valued Functions
9.6.1 Vector-Valued Functions
9.6.2 Summary
9.6.3 Exercises
9.7 Derivatives and Integrals of Vector-Valued Functions
9.7.1 The Derivative
9.7.2 Computing Derivatives
9.7.3 Tangent Lines
9.7.4 Integrating a Vector-Valued Function
9.7.5 Projectile Motion
9.7.6 Summary
9.7.7 Exercises
9.8 Arc Length and Curvature
9.8.1 Arc Length
9.8.2 Parameterizing With Respect To Arc Length
9.8.3 Curvature
9.8.4 Summary
9.8.5 Exercises
10 Derivatives of Multivariable Functions
10.1 Limits
10.1.1 Limits of Functions of Two Variables
10.1.2 Continuity
10.1.3 Summary
10.1.4 Exercises
10.2 First-Order Partial Derivatives
10.2.1 First-Order Partial Derivatives
10.2.2 Interpretations of First-Order Partial Derivatives
10.2.3 Using tables and contours to estimate partial derivatives
10.2.4 Summary
10.2.5 Exercises
10.3 Second-Order Partial Derivatives
10.3.1 Second-Order Partial Derivatives
10.3.2 Interpreting the Second-Order Partial Derivatives
10.3.3 Summary
10.3.4 Exercises
10.4 Linearization: Tangent Planes and Differentials
10.4.1 The Tangent Plane
10.4.2 Linearization
10.4.3 Differentials
10.4.4 Summary
10.4.5 Exercises
10.5 The Chain Rule
10.5.1 The Chain Rule
10.5.2 Tree Diagrams
10.5.3 Summary
10.5.4 Exercises
10.6 Directional Derivatives and the Gradient
10.6.1 Directional Derivatives
10.6.2 Computing the Directional Derivative
10.6.3 The Gradient
10.6.4 The Direction of the Gradient
10.6.5 The Length of the Gradient
10.6.6 Applications
10.6.7 Summary
10.6.8 Exercises
10.7 Optimization
10.7.1 Extrema and Critical Points
10.7.2 Classifying Critical Points: The Second Derivative Test
10.7.3 Optimization on a Restricted Domain
10.7.4 Summary
10.7.5 Exercises
10.8 Constrained Optimization: Lagrange Multipliers
10.8.1 Constrained Optimization and Lagrange Multipliers
10.8.2 Summary
10.8.3 Exercises
11 Multiple Integrals
11.1 Double Riemann Sums and Double Integrals over Rectangles
11.1.1 Double Riemann Sums over Rectangles
11.1.2 Double Riemann Sums and Double Integrals
11.1.3 Interpretation of Double Riemann Sums and Double integrals.
11.1.4 Summary
11.1.5 Exercises
11.2 Iterated Integrals
11.2.1 Iterated Integrals
11.2.2 Summary
11.2.3 Exercises
11.3 Double Integrals over General Regions
11.3.1 Double Integrals over General Regions
11.3.2 Summary
11.3.3 Exercises
11.4 Applications of Double Integrals
11.4.1 Mass
11.4.2 Area
11.4.3 Center of Mass
11.4.4 Probability
11.4.5 Summary
11.4.6 Exercises
11.5 Double Integrals in Polar Coordinates
11.5.1 Polar Coordinates
11.5.2 Integrating in Polar Coordinates
11.5.3 Summary
11.5.4 Exercises
11.6 Surfaces Defined Parametrically and Surface Area
11.6.1 Parametric Surfaces
11.6.2 The Surface Area of Parametrically Defined Surfaces
11.6.3 Summary
11.6.4 Exercises
11.7 Triple Integrals
11.7.1 Triple Riemann Sums and Triple Integrals
11.7.2 Summary
11.7.3 Exercises
11.8 Triple Integrals in Cylindrical and Spherical Coordinates
11.8.1 Cylindrical Coordinates
11.8.2 Triple Integrals in Cylindrical Coordinates
11.8.3 Spherical Coordinates
11.8.4 Triple Integrals in Spherical Coordinates
11.8.5 Summary
11.8.6 Exercises
11.9 Change of Variables
11.9.1 Change of Variables in Polar Coordinates
11.9.2 General Change of Coordinates
11.9.3 Change of Variables in a Triple Integral
11.9.4 Summary
11.9.5 Exercises
12 Vector Calculus
12.1 Vector Fields
12.1.1 Examples of Vector Fields
12.1.2 Mathematical Vector Fields
12.1.3 Plotting Vector Fields
12.1.4 Gradient Vector Fields
12.1.5 Summary
12.1.6 Exercises
12.2 The Idea of a Line Integral
12.2.1 Orientations of Curves
12.2.2 Line Integrals
12.2.3 Properties of Line Integrals
12.2.4 The Circulation of a Vector Field
12.2.5 Summary
12.2.6 Exercises
12.3 Using Parametrizations to Calculate Line Integrals
12.3.1 Parametrizations in the Definition of ∫C F⋅d r
12.3.2 Alternative Notation for Line Integrals
12.3.3 Independence of Parametrization for a Fixed Curve
12.3.4 Summary
12.3.5 Exercises
12.4 Path-Independent Vector Fields and the Fundamental Theorem of Calculus for Line Integrals
12.4.1 Path-Independent Vector Fields
12.4.2 Line Integrals Along Closed Curves
12.4.3 What other vector fields are path-independent?
12.4.4 Summary
12.4.5 Exercises
12.5 Line Integrals of Scalar Functions
12.5.1 Defining line integrals of scalar functions
12.5.2 Using Parameterizations to Calculate Scalar Line Integrals
12.5.3 Properties of Scalar Line Integrals
12.5.4 Visualizations of Scalar Line Integrals as Area Under a Curve
12.5.5 Summary
12.5.6 Exercises
12.6 The Divergence of a Vector Field
12.6.1 Definition of the Divergence of a Vector Field
12.6.2 Measuring the Change in Strength of a Vector Field
12.6.3 Summary
12.6.4 Exercises
12.7 The Curl of a Vector Field
12.7.1 Measuring the Circulation Density of Vector Field in R 2
12.7.2 Measuring Rotation in Three Dimensions
12.7.3 Circulation Density in Three Dimensions
12.7.4 Interpretation and Usage of Curl
12.7.5 Summary
12.7.6 Exercises
12.8 Green’s Theorem
12.8.1 Circulation
12.8.2 Green’s Theorem
12.8.3 What happens when vector fields are not smooth?
12.8.4 Summary
12.8.5 Exercises
12.9 Flux Integrals
12.9.1 The Idea of the Flux of a Vector Field through a Surface
12.9.2 The Details of Measuring the Flux of a Vector Field through a Surface
12.9.3 Summary
12.9.4 Exercises
12.10 Surface Integrals of Scalar Valued Functions
12.10.1 Defining surface integrals of scalar functions
12.10.2 Properties of Scalar Surface Integrals
12.10.3 Summary
12.10.4 Exercises
12.11 Stokes’ Theorem
12.11.1 Circulation in three dimensions and Stokes’ Theorem
12.11.2 Verifying and Applying Stokes’ Theorem
12.11.3 Practice with Surfaces and their Boundaries
12.11.4 Summary
12.11.5 Exercises
12.12 The Divergence Theorem
12.12.1 The Divergence Theorem
12.12.2 Summary
12.12.3 Exercises
Back Matter
Index
Section 11.4 Applications of Double Integrals
Motivating Questions
If we have a mass density function for a lamina (thin plate), how does a double integral determine the mass of the lamina?🔗 🔗
How may a double integral be used to find the area between two curves?🔗 🔗
Given a mass density function on a lamina, how can we find the lamina’s center of mass?🔗 🔗
What is a joint probability density function? How do we determine the probability of an event if we know a probability density function?🔗 🔗
🔗
So far, we have interpreted the double integral of a function f over a domain D in two different ways. First, ∬D f(x,y)d A tells us a difference of volumes — the volume the surface defined by f bounds above the x y-plane on D minus the volume the surface bounds below the x y-plane on .D. In addition, 1 A(D)∬D f(x,y)d A determines the average value of f on .D. In this section, we investigate several other applications of double integrals, using the integration process as seen in Preview Activity 11.4.1: we partition into small regions, approximate the desired quantity on each small region, then use the integral to sum these values exactly in the limit.
🔗
The following preview activity explores how a double integral can be used to determine the density of a thin plate with a mass density distribution. Recall that in single-variable calculus, we considered a similar problem and computed the mass of a one-dimensional rod with a mass-density distribution. There, as here, the key idea is that if density is constant, mass is the product of density and volume.
🔗
Preview Activity 11.4.1.
Suppose that we have a flat, thin object (called a lamina) whose density varies across the object. We can think of the density on a lamina as a measure of mass per unit area. As an example, consider a circular plate D of radius 1 cm centered at the origin whose density δ varies depending on the distance from its center so that the density in grams per square centimeter at point (x,y) is
δ(x,y)=10−2(x 2+y 2).
Suppose that we partition the plate into subrectangles ,R i j, where 1≤i≤m and ,1≤j≤n, of equal area ,Δ A, and select a point (x i j∗,y i j∗) in R i j for each i and .j. What is the meaning of the quantity ?δ(x i j∗,y i j∗)Δ A?🔗 🔗
State a double Riemann sum that provides an approximation of the mass of the plate.🔗 🔗
Explain why the double integral ∬D δ(x,y)d A tells us the exact mass of the plate. 🔗 🔗
Determine an iterated integral which, if evaluated, would give the exact mass of the plate. Do not actually evaluate the integral. (This integral is considerably easier to evaluate in polar coordinates, which we will learn more about in Section 11.5.)🔗 🔗
🔗
🔗
Subsection 11.4.1 Mass
Density is a measure of some quantity per unit area or volume. For example, we can measure the human population density of some region as the number of humans in that region divided by the area of that region. In physics, the mass density of an object is the mass of the object per unit area or volume. As suggested by Preview Activity 11.4.1, the following holds in general.
🔗
The mass of a lamina.
If δ(x,y) describes the density of a lamina defined by a planar region ,D, then the mass of D is given by the double integral .∬D δ(x,y)d A.
🔗
🔗
Activity 11.4.2.
Let D be a half-disk lamina of radius 3 in quadrants IV and I, centered at the origin as shown in Figure 11.4.1. Assume the density at point (x,y) is given by .δ(x,y)=x. Find the exact mass of the lamina.
🔗
Figure 11.4.1.A half disk lamina.
🔗
🔗
🔗
Subsection 11.4.2 Area
If we consider the situation where the mass-density distribution is constant, we can also see how a double integral may be used to determine the area of a region. Assuming that δ(x,y)=1 over a closed bounded region ,D, where the units of δ are “mass per unit of area,” it follows that ∬D 1 d A is the mass of the lamina. But since the density is constant, the numerical value of the integral is simply the area.
🔗
As the following activity demonstrates, we can also see this fact by considering a three-dimensional solid whose height is always 1.
🔗
Activity 11.4.3.
Suppose we want to find the area of the bounded region D between the curves
and y=1−x 2 and y=x−1.
🔗
A picture of this region is shown in Figure 11.4.2.
The volume of a solid with constant height is given by the area of the base times the height. Hence, we may interpret the area of the region D as the volume of a solid with base D and of uniform height 1. That is, the area of D is given by .∬D 1 d A. Write an iterated integral whose value is .∬D 1 d A. (Hint: Which order of integration might be more efficient? Why?)
Figure 11.4.2.The graphs of y=1−x 2 and .y=x−1.🔗
🔗 🔗
2. Evaluate the iterated integral from (a). What does the result tell you?🔗 🔗
🔗
🔗
We now formally state the conclusion from our earlier discussion and Activity 11.4.3.
🔗
The double integral and area.
Given a closed, bounded region D in the plane, the area of ,D, denoted ,A(D), is given by the double integral
A(D)=∬D 1 d A.
🔗
🔗
🔗
Subsection 11.4.3 Center of Mass
The center of mass of an object is a point at which the object will balance perfectly. For example, the center of mass of a circular disk of uniform density is located at its center. For any object, if we throw it through the air, it will spin around its center of mass and behave as if all the mass is located at the center of mass.
🔗
In order to understand the role that integrals play in determining the center of mass of an object with a nonuniform mass distribution, we start by finding the center of mass of a collection of N distinct point-masses in the plane.
🔗
Let ,m 1,,m 2,,…,m N be N masses located in the plane. Think of these masses as connected by rigid rods of negligible weight from some central point .(x,y). A picture with four masses is shown in Figure 11.4.3. Now imagine balancing this system by placing it on a thin pole at the point (x,y) perpendicular to the plane containing the masses. Unless the masses are perfectly balanced, the system will fall off the pole. The point (x―,y―) at which the system will balance perfectly is called the center of mass of the system. Our goal is to determine the center of mass of a system of discrete masses, then extend this to a continuous lamina.
🔗
Figure 11.4.3.A center of mass (x―,y―) of four masses.
🔗
Each mass exerts a force (called a moment) around the lines x=x― and y=y― that causes the system to tilt in the direction of the mass. These moments are dependent on the mass and the distance from the given line. Let (x 1,y 1) be the location of mass ,m 1,(x 2,y 2) the location of mass ,m 2, etc. In order to balance perfectly, the moments in the x direction and in the y direction must be in equilibrium. We determine these moments and solve the resulting system to find the equilibrium point (x―,y―) at the center of mass.
🔗
The force that mass m 1 exerts to tilt the system from the line y=y― is
m 1 g(y―−y 1),
where g is the gravitational constant. Similarly, the force mass m 2 exerts to tilt the system from the line y=y― is
m 2 g(y―−y 2).
🔗
In general, the force that mass m k exerts to tilt the system from the line y=y― is
m k g(y―−y k).
🔗
For the system to balance, we need the forces to sum to 0, so that
∑k=1 N m k g(y―−y k)=0.
🔗
Solving for ,y―, we find that
y―=∑k=1 N m k y k∑k=1 N m k.
🔗
A similar argument shows that
x―=∑k=1 N m k x k∑k=1 N m k.
🔗
The value M x=∑k=1 N m k y k is called the total moment with respect to the x-axis; M y=∑k=1 N m k x k is the total moment with respect to the y-axis. Hence, the respective quotients of the moments to the total mass, ,M, determines the center of mass of a point-mass system:
(x―,y―)=(M y M,M x M).
🔗
Now, suppose that rather than a point-mass system, we have a continuous lamina with a variable mass-density .δ(x,y). We may estimate its center of mass by partitioning the lamina into m n subrectangles of equal area ,Δ A, and treating the resulting partitioned lamina as a point-mass system. In particular, we select a point (x i j∗,y i j∗) in the i j th subrectangle, and observe that the quanity
δ(x i j∗,y i j∗)Δ A
is density times area, so δ(x i j∗,y i j∗)Δ A approximates the mass of the small portion of the lamina determined by the subrectangle .R i j.
🔗
We now treat δ(x i j∗,y i j∗)Δ A as a point mass at the point .(x i j∗,y i j∗). The coordinates (x―,y―) of the center of mass of these m n point masses are thus given by
and x―=∑j=1 n∑i=1 m x i j∗δ(x i j∗,y i j∗)Δ A∑j=1 n∑i=1 m δ(x i j∗,y i j∗)Δ A and y―=∑j=1 n∑i=1 m y i j∗δ(x i j∗,y i j∗)Δ A∑j=1 n∑i=1 m δ(x i j∗,y i j∗)Δ A.
🔗
If we take the limit as m and n go to infinity, we obtain the exact center of mass (x―,y―) of the continuous lamina.
🔗
The center of mass of a lamina.
The coordinates (x―,y―) of the center of mass of a lamina D with density δ=δ(x,y) are given by
and x―=∬D x δ(x,y)d A∬D δ(x,y)d A and y―=∬D y δ(x,y)d A∬D δ(x,y)d A.
🔗
🔗
The center of mass of a lamina can then be thought of as a weighted average of all of the points in the lamina with the weights given by the density at each point. The centroid of a lamina is the just the average of all of the points in the lamina, or the center of mass if the density at each point is 1.
🔗
The numerators of x― and y― are called the respective moments of the lamina about the coordinate axes. Thus, the moment of a lamina D with density δ=δ(x,y) about the y-axis is
M y=∬D x δ(x,y)d A
and the moment of D about the x-axis is
M x=∬D y δ(x,y)d A.
🔗
If M is the mass of the lamina, it follows that the center of mass is
.(x―,y―)=(M y M,M x M).
🔗
Activity 11.4.4.
In this activity we determine integrals that represent the center of mass of a lamina D described by the triangular region bounded by the x-axis and the lines x=1 and y=2 x in the first quadrant if the density at point (x,y) is .δ(x,y)=6 x+6 y+6. A picture of the lamina is shown in Figure 11.4.4.
🔗
Figure 11.4.4.The lamina bounded by the x-axis and the lines x=1 and y=2 x in the first quadrant.
🔗
Set up an iterated integral that represents the mass of the lamina.🔗 🔗
Assume the mass of the lamina is 14. Set up two iterated integrals that represent the coordinates of the center of mass of the lamina.🔗 🔗
🔗
🔗
Subsection 11.4.4 Probability
Calculating probabilities is a very important application of integration in the physical, social, and life sciences. To understand the basics, consider the game of darts in which a player throws a dart at a board and tries to hit a particular target. Let us suppose that a dart board is in the form of a disk D with radius 10 inches. If we assume that a player throws a dart at random, and is not aiming at any particular point, then it is equally probable that the dart will strike any single point on the board. For instance, the probability that the dart will strike a particular 1 square inch region is ,1 100 π, or the ratio of the area of the desired target to the total area of D (assuming that the dart thrower always hits the board itself at some point). Similarly, the probability that the dart strikes a point in the disk D 3 of radius 3 inches is given by the area of D 3 divided by the area of .D. In other words, the probability that the dart strikes the disk D 3 is
9 π 100 π=∬D 3 1 100 π d A.
🔗
The integrand, ,1 100 π, may be thought of as a distribution function, describing how the dart strikes are distributed across the board. In this case the distribution function is constant since we are assuming a uniform distribution, but we can easily envision situations where the distribution function varies. For example, if the player is fairly good and is aiming for the bulls eye (the center of D), then the distribution function f could be skewed toward the center, say
f(x,y)=K e−(x 2+y 2)
for some constant positive .K. If we assume that the player is consistent enough so that the dart always strikes the board, then the probability that the dart strikes the board somewhere is 1, and the distribution function f will have to satisfy
1
This makes ,K=1 π(1−e−100), which you can check.
∬D f(x,y)d A=1.
🔗
For such a function ,f, the probability that the dart strikes in the disk D 1 of radius 1 would be
∬D 1 f(x,y)d A.
🔗
Indeed, the probability that the dart strikes in any region R that lies within D is given by
∬R f(x,y)d A.
🔗
The preceding discussion highlights the general idea behind calculating probabilities. We assume we have a joint probability density function,f, a function of two independent variables x and y defined on a domain D that satisfies the conditions
f(x,y)≥0 for all x and y in ,D,🔗 🔗
the probability that x is between some values a and b while y is between some values c and d is given by ∫a b∫c d f(x,y)d y d x, 🔗 🔗
The probability that the point (x,y) is in D is 1, that is (11.4.1)(11.4.1)∬D f(x,y)d A=1. 🔗 🔗
🔗
Note that it is possible that D could be an infinite region and the limits on the integral in Equation(11.4.1) could be infinite. When we have such a probability density function ,f=f(x,y), the probability that the point (x,y) is in some region R contained in the domain D (the notation we use here is “P((x,y)∈R)”) is determined by
P((x,y)∈R)=∬R f(x,y)d A.
🔗
Activity 11.4.5.
A firm manufactures smoke detectors. Two components for the detectors come from different suppliers — one in Michigan and one in Ohio. The company studies these components for their reliability and their data suggests that if x is the life span (in years) of a randomly chosen component from the Michigan supplier and y the life span (in years) of a randomly chosen component from the Ohio supplier, then the joint probability density function f might be given by
f(x,y)=e−x e−y.
Theoretically, the components might last forever, so the domain D of the function f is the set D of all (x,y) such that x≥0 and .y≥0. To show that f is a probability density function on D we need to demonstrate that ∫∫D f(x,y)d A=1, or that ∫0∞∫0∞f(x,y)d y d x=1. Use your knowledge of improper integrals to verify that f is indeed a probability density function. 🔗 🔗
Assume that the smoke detector fails only if both of the supplied components fail. To determine the probability that a randomly selected detector will fail within one year, we will need to determine the probability that the life span of each component is between 0 and 1 years. Set up an appropriate iterated integral, and evaluate the integral to determine the probability.🔗 🔗
What is the probability that a randomly chosen smoke detector will fail between years 3 and 7?🔗 🔗
Suppose that the manufacturer determines that one of the components is more likely to fail than the other, and hence conjectures that the probability density function is instead f(x,y)=K e−x e−2 y. What is the value of ?K?🔗 🔗
🔗
🔗
🔗
Subsection 11.4.5 Summary
The mass of a lamina D with a mass density function δ=δ(x,y) is ∬D δ(x,y)d A.🔗 🔗
The area of a region D in the plane has the same numerical value as the volume of a solid of uniform height 1 and base ,D, so the area of D is given by ∬D 1 d A.🔗 🔗
The center of mass, ,(x―,y―), of a continuous lamina with a variable density δ(x,y) is given by and x―=∬D x δ(x,y)d A∬D δ(x,y)d A and y―=∬D y δ(x,y)d A∬D δ(x,y)d A. 🔗 🔗
Given a joint probability density function f is a function of two independent variables x and y defined on a domain ,D, if R is some subregion of ,D, then the probability that (x,y) is in R is given by ∬R f(x,y)d A. 🔗 🔗
🔗
Exercises 11.4.6 Exercises
1.
Activate
The masses m i are located at the points .P i. Find the center of mass of the system.
🔗
,m 1=1,,m 2=2,.m 3=3.
🔗
,P 1=(9,−4),,P 2=(−6,−8),.P 3=(5,2).
🔗
x¯=
🔗
y¯=
🔗
🔗
2.
Activate
Find the centroid (x¯,y¯) of the triangle with vertices at ,(0,0),,(7,0), and .(0,4).
🔗
x¯=
🔗
y¯=
🔗
🔗
3.
Activate
Find the mass of the rectangular region ,0≤x≤2,0≤y≤2 with density function .ρ(x,y)=2−y.
🔗
🔗
4.
Activate
Find the mass of the triangular region with vertices (0, 0), (4, 0), and (0, 1), with density function .ρ(x,y)=x 2+y 2.
🔗
🔗
5.
Activate
A lamina occupies the region inside the circle x 2+y 2=12 y but outside the circle .x 2+y 2=36. The density at each point is inversely proportional to its distance from the orgin.
🔗
Where is the center of mass?
🔗
( , )
🔗
🔗
6.
Activate
A sprinkler distributes water in a circular pattern, supplying water to a depth of e−r feet per hour at a distance of r feet from the sprinkler.
🔗
A. What is the total amount of water supplied per hour inside of a circle of radius 9?
🔗
f t 3 per hour
🔗
B. What is the total amount of water that goes through the sprinkler per hour?
🔗
f t 3 per hour
🔗
🔗
7.
Activate
Let p be the joint density function such that p(x,y)=1 9 x y in ,R, the rectangle ,0≤x≤6,0≤y≤1, and p(x,y)=0 outside .R. Find the fraction of the population satisfying the constraint x+y≤7
🔗
fraction =
🔗
🔗
8.
Activate
A lamp has two bulbs, each of a type with an average lifetime of 2 hours. The probability density function for the lifetime of a bulb is .f(t)=1 2 e−t/2,t≥0.
🔗
What is the probability that both of the bulbs will fail within 5 hours?
🔗
🔗
9.
Activate
For the following two functions ,p(x,y), check whether p is a joint density function. Assume p(x,y)=0 outside the region .R.
🔗
(a),p(x,y)=1, where R is .0≤x≤1,2≤y≤2.5.
🔗
p(x,y)
is a joint density function🔗 🔗
is not a joint density function🔗 🔗
🔗
(b),p(x,y)=2, where R is .2≤x≤2.5,−1≤y≤3.
🔗
p(x,y)
is a joint density function🔗 🔗
is not a joint density function🔗 🔗
🔗
Then, for the region R given by ,2≤x≤3,0≤y≤2, what constant function p(x,y) is a joint density function?
🔗
p(x,y)=
🔗
🔗
10.
Activate
Let x and y have joint density function
for otherwise.p(x,y)={2 3(x+2 y)for 0≤x≤1,0≤y≤1,0 otherwise.
Find the probability that
🔗
(a):x>1/3:
🔗
probability =
🔗
(b):x<1 3+y:
🔗
probability =
🔗
🔗
11.
A triangular plate is bounded by the graphs of the equations ,y=2 x,,y=4 x, and .y=4. The plate’s density at (x,y) is given by ,δ(x,y)=4 x y 2+1, measured in grams per square centimeter (and x and y are measured in centimeters).
Set up an iterated integral whose value is the mass of the plate. Include a labeled sketch of the region of integration. Why did you choose the order of integration you did?🔗 🔗
Determine the mass of the plate.🔗 🔗
Determine the exact center of mass of the plate. Draw and label the point you find on your sketch from (a).🔗 🔗
What is the average density of the plate? Include units on your answer.🔗 🔗
🔗
🔗
12.
Let D be a half-disk lamina of radius 3 in quadrants IV and I, centered at the origin as in Activity 11.4.2. Assume the density at point (x,y) is equal to .x.
Before doing any calculations, what do you expect the y-coordinate of the center of mass to be? Why?🔗 🔗
Set up iterated integral expressions which, if evaluated, will determine the exact center of mass of the lamina.🔗 🔗
Use appropriate technology to evaluate the integrals to find the center of mass numerically.🔗 🔗
🔗
🔗
13.
Let x denote the time (in minutes) that a person spends waiting in a checkout line at a grocery store and y the time (in minutes) that it takes to check out. Suppose the joint probability density for x and y is
f(x,y)=1 8 e−x/4−y/2.
What is the exact probability that a person spends between 0 to 5 minutes waiting in line, and then 0 to 5 minutes waiting to check out?🔗 🔗
Set up, but do not evaluate, an iterated integral whose value determines the exact probability that a person spends at most 10 minutes total both waiting in line and checking out at this grocery store.🔗 🔗
Set up, but do not evaluate, an iterated integral expression whose value determines the exact probability that a person spends at least 10 minutes total both waiting in line and checking out, but not more than 20 minutes.🔗 🔗
🔗
🔗
🔗
🔗
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13738 | https://www.khanacademy.org/science/hs-chemistry/x2613d8165d88df5e:nuclear-chemistry-hs/x2613d8165d88df5e:half-life-and-radiometric-dating/v/half-life-radiometric-dating | Use of cookies
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13739 | https://arxiv.org/pdf/1901.00582 | ALTERNATING KNOTS
WILLIAM W. MENASCO
Abstract. This is a short expository article on alternating knots and is to appear in the Concise Encyclopedia of Knot Theory.
Introduction
Figure 1. P.G. Tait’s first knot table where he lists all knot types up to 7 crossings. (From reference , courtesy of J. Hoste, M. Thistlethwaite and J. Weeks.) A knot K ⊂ S3 is alternating if it has a regular planar diagram DK ⊂ P(∼=
S2) ⊂ S3 such that, when traveling around K , the crossings alternate, over-under-over-under, all the way along K in DK . Figure 1 show the first 15 knot types in P. G. Tait’s earliest table and each diagram exhibits this alternating pattern. This simple definition is very unsatisfying. A knot is alternating if we can draw it as an alternating diagram? There is no mention of any geometric structure. Dissatisfied with this characterization of an alternating knot, Ralph Fox (1913-1973) asked: ”What is an alternating knot?” black
white
black
white
Figure 2. Going from a black to white region near a crossing.
1
arXiv:1901.00582v1 [math.GT] 3 Jan 2019 2WILLIAM W. MENASCO
Let’s make an initial attempt to address this dissatisfaction by giving a different characterization of an alternating diagram that is immediate from the over-under-over-under characterization. As with all regular planar diagrams of knots in S3, the regions of an alternating diagram can be colored in a checkerboard fashion. Thus, at each crossing (see figure 2) we will have “two” white regions and “two” black regions coming together with similarly colored regions being kitty-corner to each other. (I
Figure 3. Black and white regions. will explain my use of quotation marks momentarily.) In the local picture of figure 2, we will make the convention that when standing in a black region and walking to a white region by stepping over the under-strand of the knot the over-strand is to our left. (Notice this is independent of which side of P we stand.) Now the feature for an alternating diagram is, if one crossing of a black region satisfies this local scheme then every crossing of that black region satisfies this scheme. Since both black regions at a crossing satisfy our scheme then this scheme is transmitted to every black region of an alternating diagram. To summarize, a diagram of a knot is alternating if we can checkerboard color the regions such that when we stand in any black region, we can walk to an adjacent white region by stepping over the under-strand of a crossing while having the over-strand of that crossing positioned to our left. From the perspective of walking from a white regions to a black regions, the over-crossing is to our right. Again, one should observe that this coloring scheme is consistent independent of which side of P one stands. Coming back now to our previous use of quotation marks—“two” white regions and “two” black regions—we must allow for the possibility that “two” kitty-corner regions may be the same region as shown in figure 4. This possibility forces the occurrence of a nugatory crossing in the diagram of DK . Such a crossing can be eliminated by a π-rotation of K′ for half the diagram. Observe that such a π-rotation K K’ L K
Figure 4. Eliminating a nugatory crossing. will take an alternating diagram to an alternating diagram with one less crossing. So the elimination of nugatory crossings must terminate and we call a diagram with no nugatory crossings reduced . Thus, our study of alternating diagrams is legitimately restricted to ones that are reduced. And, let’s add one further restriction, that a diagram DK be connected, since one that is not connected is obviously representative of a split link. This requirement implies that all regions of a diagram DK are discs ALTERNATING KNOTS 3
With all of the above introductory material in place, and restricting to reduced connected alternating knot diagrams, we now give in to the urge that anyone who has drawn a M¨ obius band will have. At each crossing we place a half-twisted white band connecting the two white regions; similarly, we place a half-twisted black band connecting the two black regions. This construction gives us a white surface, ωK , and a black surface, βK , each one having our knot K as its boundary—they are spanning surfaces of K. The white and black bands of each crossing intersect in a single arc that has one endpoint on the under-crossing strand and its other endpoint on the over-crossing strand. Thus, we can reinterpret a checkerboard DK diagram to be that of a black and white surface with each crossing concealing an intersection arc since our eye’s viewpoint would be one where we were looking straight down each arc. Here is where I really wish to start our story of alternating knots. (The story can be expanded to alternating links, but we will leave that for another time.) At the center of this story lies three mysteries, the three conjectures of the Scottish physicist Peter G. Tait (1831-1901). As with all good conjectures, Tait’s three point to deeper mathematics. Alternating knots have proven particularly well behaved with respect to tabulation, classification, and computing and topologically interpreting algebraic invariants. Moreover, for hyperbolic alternating knots their alternating diagrams are more closely tied to their hyperbolic geometric structure than other knots. As a collection, alternating knots have supplied researchers with a ready population for experimentation and testing conjectures. And, our two colored surfaces, βK and
ωK , are rightfully thought of as the heroes the story for they are at the center of establishing the validity of the Tait conjectures and, finally, giving a satisfactory answer to Fox’s question.
The Tait conjectures
The early efforts of knot tabulation came initially from Tait. He was motivated by Lord Kelvin’s program for understanding the different chemical elements as differ-ent knotted vortices in ether. Without the aid of any theorems from topology, Tait published in 1876 his first knot table which containing the 15 knot types through seven crossings in figure 1. Specifically, Tait enumerated all possible diagrams up to a seven crossings and then grouped together those diagrams that represented the same knot type. For example, the right-handed trefoil has two diagrams, one as a closed 2-braid and one as a 2-bridge knot. Similarly, the figure-eight two diagrams, one coming from a closed 3-braid and one coming from a 2-bridge presentation. His grouping of these 15 knot types is consistent with todays modern tables. However, his grouping did contained more than these 15 since he did not possess the notion of prime/composite knots—he included the composite sums of the trefoil with itself and the figure-eight. Tait also experimented with the assigning of orientations to the knot diagrams— giving a defined direction to the knot along which to travel—which induced a hand-edness to each crossing. Thus, for right-handed crossings one could assign a +1 parity and a −1 parity for left-handed crossings. (See figure 5.) The sum of these parities 4 WILLIAM W. MENASCO left-handed or -1 right-handed or +1
Figure 5. Left and right hand crossings. was the writhe of a diagram. For example, independent of orientation the trefoil in Tait’s table will have three right-handed crossings yielding a writhe of +3, whereas the figure-eight will always have two right-handed and two left-handed crossings for a writhe of 0. Tait also observed that one could reduce crossing number by the elimination of nugatory crossings. Finally, Tait discovered an operation on an alternating diagram— a flype —that preserved the number of crossings, the writhe and, more importantly, the knot type. (See figure 6.) T
T
Figure 6. A flype moves the crossing from right to left. From the diagrams and their groupings into just 15 knot types, Tait had amassed a sizable about of data. And, from this data Tait proposed a set of conjectures:
The Tait Conjectures
(T1) A reduced alternating diagram has minimal crossing number. (T2) Any two reduced alternating diagrams of the same knot type have the same writhe. (T3) Any two reduced alternating diagrams of a given knot type are related by a sequence of flypes. One should observe that T1 and T3 imply T2. The solutions to these three conjectures would have to wait until after Vaughan Jones’ work on his knot polynomial in 1984.
Surfaces and alternating knots
During the 20 th Century, work on understanding the topology and geometry of alternating knots remained active. But, for our story line we start in the 1980’s and describe how our colored surfaces, βK and ωK control the behavior of closed embedded surfaces in an alternating knot complement S3 − K. Any embedded closed surface, Σ, in S3 is necessarily orientable and the ones that capture meaningful information about the topology of a knot complement are incompressible . That is, every simple closed curve (s.c.c.) c ⊂ Σ( ⊂ S3 − K) which bounds an embedded disc in S3 − K,also bounds a sub-disc of Σ. To this end, it is worth observing from the start that ALTERNATING KNOTS 5
S3 − (βK ∪ ωK ) is the disjoint union of two open 3-balls which we denote by B3
N
(N
for north) and B3
S
(S for south). Given an incompressible surface Σ ⊂ S3 −K our first task is to put Σ into normal position with respect to βK and ωK . Using some basic general position arguments for incompressibility, this means that we can assume: (1) βK ∩ Σ and ωK ∩ Σ are collections of s.c.c.’s; (2) any curve from either collection cannot be totally contained in a single black or white region used to construct βK or ωK —it must intersect some of the half-twisted bands of our colored surfaces—and; (3) B3
N
∩ Σ and B3
S
∩ Σ are collections of open 2-discs which we will call domes .Of particular interest is the the behavior of the boundary curves of the domes. Let δ ⊂ B3
N
∩ Σ be a dome and consider the curve ∂(¯δ) ⊂ βK ∪ ωK . When viewed from inside B3
N
, the s.c.c. ∂(¯δ) will be seen as a union of arcs—an arc in a white region adjoined to an arc in a black region adjoined to an arc in a white region, etc.—where the adjoining of two consecutive arcs occurs within an intersection arc of
βK ∩ ωK . Thus, from the viewpoint of B3
N
∂(¯δ) respects our earlier scheme of when traveling from white to black (resp. black to white) regions having the over-strand to the left (resp. right). One additional normal position condition we can require by general position arguments is that any ∂(¯δ) intersects any intersection arc of βK ∩ ωK
at most once. Now, observe that if we have a curve ∂(¯δ) occurring on one side of an over-strand at a crossing then we must see a curve ∂( ¯δ′) on the another side of the same over-strand. Finally, observe that if ∂(¯δ) = ∂( ¯δ′) then Σ will contain a s.c.c. that is a meridian of K. Once placed in the above described normal position, it is
Figure 7. The view of ∂(¯δ) inside B3
N
. Traveling along ∂(¯δ), we find a crossing to left/right then right/left. readily observed by a nesting argument that Σ must contain a meridian curve of K.If we walk along the boundary curve of any dome, starting in a white region of ωK
and passing into a black region of βK , our scheme dictates that the over-strand of a crossing be to our left; then traveling through a black region into a white region our scheme says the over-strand of a crossing is to our right. (See figure 7.) After strolling along all of our boundary curve, if we have not walked on the two sides of the some crossing then there are at least two additional curves that are boundaries curves of differing domes, one to the left side of our first crossing and one to the right side of our second crossing. Walking along either one of these boundary curves and iterating this procedure we will either find a new boundary of a dome that we have not strolled along or we will find a boundary that passes on both sides of the over-strand of a crossing, thus yielding a meridian of K. By the compactness of Σ there are only finitely many domes in B3
N
so the latter must occur. 6 WILLIAM W. MENASCO
The implication of an incompressible Σ always containing a meridian curve are significant. Once a meridian curve has been found we can trade Σ in for a new surface punctured by K, Σ ′, by compressing Σ along a once puncture disc that the meridian curve bounds. (See figure 8.) Repeating this meridian-surgery whenever a new meridian curve is found, the study of closed incompressible surfaces in the complement of alternating knots becomes the study of incompressible meridionally or
pairwise incompressible surfaces—if a curve in Σ bounds a disc (resp. once punctured disc) in S3 − K then it bounds a disc (resp. once punctured disc) in Σ. Finally, the L
Figure 8. Pairwise compressing Σ along a meridian disc. above normal position and nesting argument can be enhanced so as to handle the intersection of an incompressible pairwise incompressible surface Σ. Specifically, if
K is an alternating composite knot then there exists a twice punctured 2-sphere the intersects βK and ωK in single arcs whose union forms a circle that can be thought of as illustrating the composite nature of DK .Thus, we have the result that an alternating knot is a composite knot if and only if its diagram is also composite. When combined with William P. Thurston’s characterization of hyperbolic knots in S3, an immediate consequence is that all alternating knots coming from diagrams that are neither composite nor (2 , q )-torus knots are hyperbolic knots.
Hyperbolic geometry and alternating knots
In 1975 Robert Riley produced the first example of a hyperbolic knot, the figure eight, by showing that its fundamental group had a faithful discrete representation into P SL (2 , C), the group of orientation preserving isometries of hyperbolic 3-space. Inspired by Riley’s pioneering work (see page 360 of reference 2), Thurston proved that all knots that are not torus knots and not satellite knots are hyperbolic—their complement has a complete hyperbolic structure. Recalling that a satellite knot is one where there is an incompressible torus that is not a peripheral torus, we now apply the fact that for alternating knots such a torus has a meridian curve. Then pairwise compressing such a torus will produce an annulus illustrating that the knot is a composite knot. As mentined previously, we then conclude that all prime alternating knots that are knot (2 , q )-torus knots are hyperbolic. Similar to Riley, Thurston was able to structure the hyperbolic structures for knot complements. However, he used a geometric approach to see their hyperbolic structures. Specifically, the complement of every hyperbolic knot can be decomposed in a canonical way into ideal tetrahedra, that is, convex tetrahedra in hyperbolic 3-space such that all the vertices are lying on the sphere at infinity. The most famous example of such a tetrahedron decomposition is Thurston’s decomposition of the figure eight complement into two regular ideal tetrahedra. of ALTERNATING KNOTS 7
our previous two open 3-balls, B3
N
and B3
S
. Referring to figure 9, if we are viewing the diagram from inside B3
N
and we try to look into B3
N
, our view is blocked by the non-opaque black and white disc-regions of βK and ωK . Moreover, when we look black and white regions are adjoined black and white regions are adjoined
Figure 9. From the viewpoint of B3
N
, black and white regions are adjoined along an (oriented) edge). at an arc of βK ∩ ωK , from the perspective of B3
N
it does not appear that there are two black (white) regions adjoined be a half-twisted band. Instead it appears that the black regions are adjoined to the white regions in a manner consistent with our original scheme of going from a white region and black region and having the over-strand to our left. This adjoining scheme for assembling the boundary of B3
N
,and similarly B3
S
from the disc-regions of our colored surfaces makes sense once one realizes that one’s vision in B3
N
gives us information for only the over-strand at each crossing since we see it in its entirety. To obtain the information for the under-strand of a crossing—to see it in its entirety—we need to view it in B3
S
where it appears as an over-strand of a crossing. In figure 10 we show the inside appearance of the A BCDABDCABDCAA
bigons be byegones bigons be byegones
Figure 10. The upper left (resp. right) configurations are the view of
SN (resp. SS ) from inside B3
N
(resp. B3
S
). 2-sphere boundaries, SN = ∂B 3
N
and SS = ∂B 3
S
, coming from this described scheme for assembling of the disc-regions of our colored surface. Notice that SN and SS each has four vertices—one for each crossing. When the identification of commonly labeled disc-regions of SN and SS is made these 8 vertices will become one vertex and the resulting space will correspond to the space obtained by taking the complement an open tubular neighborhood of the figure eight and coning its peripheral torus to a point. Thus, when we delete this single vertex we have the knot complement. When we place B3
N
and B3
S
in hyperbolic 3-space, H3, so that the vertices of SN and SS are in S∞, the sphere at infinity—they are ideal polyhedra—we have achieved this vertex 8 WILLIAM W. MENASCO
deletion. A complete hyperbolic structure for the knot complement will come from a tessellation of H3 by isometric copies of our two ideal polyhedra. Now, one feature of an ideal polyhedron is that each edge is the unique geodesic line between two points in S∞. Currently since some edges are boundaries of bigon disc-regions in SN and
SS we need to collapse them to achieve the needed unique edge—we let bigons be byegones. For the figure eight, once the four bigon disc-regions are collapsed on SN
and SS , B3
N
and B3
S
become Thurston’s two tetrahedra decomposition. The above assembly scheme for producing a ideal polyhedron works for any alternating knot that is not a (2 , q )-torus or composite knot with only one caveat. If a disc-region is not a triangle, we need to insert additional edges until it is a union of triangles.
Establishing conjecture T1
In 1985 Vaughan Jones announce the discovery of his polynomial, VK (t), a Lau-rent polynomial that is an invariant of oriented knots (and links). His discovery emerged from a study of finite dimensional von Neumann algebras. A number of topologists subsequently reframed the Jones polynomial into more knot theory user friendly settings. Then in 1987 using these reframings of the Jones polynomial, Louis Kauffman, Kunio Murasugi, and Morwen B. Thistlethwaite each independently pub-lished proofs of the conjecture T1, a reduced alternating diagram is a minimal crossing diagram . Additionally, both Mursugi and Thiethlewaite showed that the writhe was an invariant of a reduced alternating diagram, establishing conjecture T2. Kauffman’s approach using his bracket invariant and state model equation is particularly accessible. Moreover, his approach illustrates the importance of our two colored surfaces. We consider an n-crossing knot diagram, D, where we have given each crossing a label {1, · · · , n }. A state of D is a choice of smoothly resolving every crossing. The ith crossing can be resolve either positively or negatively as shown in figure 11. s(i)=-1 s(i)=+1
Figure 11. Kauffman’s states of D.Thus, a state of D also corresponds to a function s : {1, · · · , n } → {− 1, +1 }
where s(i) = +1 (resp . −1) if we smooth the ith crossing in a positive (resp. negative) manner. Applying a state s to D we obtain a new diagram, sD , that is just a collection of disjoint simple closed curves—there are no crossings. We will use |sD | to mean the number of curves in the collection. Of immediate interest is when D is a reduced alternating diagram and the state is constant, s+ ≡ +1. A constant +1 state can be thought of as cutting every half-twisted band of ωK . Notice then |s+D| is just equal to the number of disc-regions of
ωK . Similarly, for |s−D| is a count of the number of disc-regions of βK . Moreover, ALTERNATING KNOTS 9
using the classical Euler characteristic relation for the sphere, (vertices − edges + faces) = 2, we can derive the equality |s+D| + |s−D| = n + 2. We now come to the Kauffman bracket of D which is a Laurent polynomial, 〈D〉,with an indeterminate variable A. Initially Kauffman axiomatically define his bracket in terms of skein relations. He then showed it was an invariant of regular isotopies ,isotopies that corresponded to a sequence of only Reidemeister type II and III moves. Later, Kauffman derived the following state sum model for his bracket (1) 〈D〉 = ∑
s
(A∑ni=1 s(i)(−A−2 − A2)|sD |− 1).
Although the Kauffman bracket is only a regular isotopy invariant, when we insert a term that takes the writhe, w(D), into account we do obtain an oriented knot invariant, the Kauffman polynomial :(2) FK (A) = ( −A)−3w(DK )〈DK 〉.
As an aside we mention that the Jones polynomial can be obtained from the Kauffman polynomial by a variable substantiation, FK (t14 ) = VK (t). Now since FK (A) is an oriented knot invariant, the difference between the high-est exponent power, M (FK ), and lowest exponent power, m(FK ), of A in FK (A) is also an invariant. But, since the ( −A)−3w(DK ) factor in equation (2) is common to both M (FK ) and m(FK ) we see that M (FK ) − m(FK ) = M (〈DK 〉) − m(〈DK 〉), the difference of the highest and lowest power of A in 〈DK 〉. Now referring back to equa-tion (1) with a little arguing one can obtain that M (〈DK 〉) = n + 2 |s+D| − 2 and
m(〈DK 〉) = −n − 2|s−D| + 2. Then:
M (〈DK 〉) − m(〈DK 〉) = 2 n + 2( |s+D| + |s−D|) − 2.
Finally, recalling |s+D| + |s−D| = n + 2 we obtain M (〈DK 〉) − m(〈DK 〉) = 4 n.A variation of the above line of reasoning can establish that for an n-crossing non-alternating prime diagram, D, we have M (〈D〉) − m(〈D〉) < 4n. But, since
M (〈D〉) − m(〈D〉) is an invariant we conclude that a reduced alternating diagram has minimal number of crossings.
Establishing conjectures T2 & T3
Following the successful verification of conjecture T1 in 1987, Thistlethwaite and William W Menasco announced in 1991 a proof verifying conjecture T3, the Tait fly-ping conjecture. As mentioned before T1 and T3 together imply T2. Their proof of T3 had two components. First, it utilized the new understanding of the interplay be-tween the algebraic invariants coming out of the Jones polynomial revolution and the crossing number of knot diagrams. Second, it utilized enhancements of the elemen-tary geometric technology (previously described) in order to analyze the intersection pattern of a closed incompressible surfaces with the colored surfaces, β ∪ ω, of an alternating diagram. Two key lemmas of this enhanced geometric technology are worth mentioning. First, for a reduced alternating diagram the associated colored surfaces, β and ω, are 10 WILLIAM W. MENASCO
themselves incompressible. Second, for a reduced alternating diagram neither β nor
ω contain a meridional simple closed curve. The overall strategy of the proof is fairly straight forward. Given two reduced alternating diagrams, D1 and D2, representing the same knot type, take the associated two sets of incompressible colored surfaces, β1 ∪ ω1 and β2 ∪ ω2, and consider how they will intersect each other. Using nesting arguments that are similar in spirit to the one for closed surfaces, one can establish that intersections are arcs occurring away from β1 ∩ ω1 and β2 ∩ ω2—the portion of our colored surfaces that come from the half-twisted bands. The argument proceeds by trying to “line up” crossings of
D1 with crossings of D2. An extensive analysis of these arc intersections then leads to the conclusion that if D1 and D2 are not isotopic diagrams, there exists a flype. After performing a sufficient number of the flypes we have more crossings lining up. When all the crossings of D1 line up with those of D2 we have isotopic diagrams. With the establishment of T3 we have the structure for a complete classification of alternating knots.
What is an alternating knot?
In 2017 independently Joshua E. Greene and Joshua Howie published differing answers to Fox’s fundamental question. Howie’s answer starts in the setting of knots in S3 and supposes that such a knot, K, admits two spanning surfaces, Σ and Σ ′.Neither surface need be orientable. Let i(∂Σ, ∂ Σ′) be the intersection number of the the two curves obtained by intersecting Σ ∪ Σ′ with a peripheral torus. Then Howie proves that K is an n-crossing alternating knot in S3 if and only if the following Euler characteristic equation is satisfied:
χ(Σ) + χ(Σ ′) + 12i(∂Σ, ∂ Σ′) = 2 .
The “only if” direction of this statement makes sense once one realizes that 12 i(∂β, ∂ω ) =
n, since each arc of β ∩ ω accounts for two intersections on the peripheral torus, and
χ(Σ) + χ(Σ ′) = ( |s+D| − n) + ( |s−D| − n). Greene’s answer starts in a more general setting, knots in a Z2-homology 3-spheres, M 3. He also supposes that one has two spanning surfaces, Σ and Σ ′, of a knot, K ⊂ M 3. One then considers the Gordon-Litherland pairing form of the 1st -homology for each surface:
FΣ : H1(Σ) × H1(Σ) → Z and FΣ′ : H1(Σ ′) × H1(Σ ′) → Z.
If FΣ is a positive definite form and FΣ′ is a negative definite form then Greene’s result states that M 3 = S3 and K is an alternating knot. Moreover, Σ and Σ ′ are our colored surfaces—our two heroes—for an alternating diagram, DK .ALTERNATING KNOTS 11
References and further readings Introduction.
P. G. Tait, On knots I, II, III, Scientific Papers, Cambridge University Press, 1898-1900. Including: Trans. R. Soc. Edin., 28, 1877, 35-79. Reprinted by Amphion Press, Washington D.C., 1993. 2. T. P Kirkman, The enumeration, description and construction of knots of fewer than ten crossings, Trans. Roy. Soc. Edinburgh 32 (1885), 281-309. 3. T. P Kirkman, The 364 unifilar knots of ten crossings enumerated and defined, Trans. Roy. Soc. Edinburgh 32 (1885), 483-506. 4. C. N. Little, On knots, with a census of order ten, Trans. Connecticut Acad. Sci. 18 (1885), 374-378. 5. C. N. Little, Non alternate knots of orders eight and nine, Trans. Roy. Soc. Edinburgh 35 (1889), 663-664. 6. Jim Hoste, Morwen B. Thistlethwaite and Jeffrey Weeks, The first 1,701,936 knots, Math. Intelligencer 20 (4) (1998), 33-48
The Tait conjectures & Surfaces and alternating knots.
William W. Menasco, Closed incompressible surfaces in alternating knot and link complements, Topology 23 (1) (1984), 37-44. 2. William W. Menasco, Determining incompressibility of surfaces in alternating knot and link complements. Pacific J. Math. 117 (1985), no. 2, 353-370. 3. William W. Menasco and M. B Thistlethwaite, Surfaces with boundary in alternating knot exteriors. J. Reine Angew. Math. 426 (1992), 47-65.
Hyperbolic geometry and alternating knots.
Robert Riley, Discrete parabolic representations of link groups, Mathematika, 22 (2) (1975), pp. 141-150. 2. William P. Thurston, Three dimensional manifolds, Kleinian groups and hy-perbolic geometry, Bull. Amer. Math. Soc. (NS), 6 (1982), pp. 357-381 3. William W. Menasco, Polyhedra representation of link complements. Low-dimensional topology (San Francisco, Calif., 1981), 305?325, Contemp. Math., 20, Amer. Math. Soc., Providence, RI, 1983. 4. Marc Lackenby, The volume of hyperbolic alternating link complements. With an appendix by Ian Agol and Dylan Thurston. Proc. London Math. Soc. (3) 88 (2004), no. 1, 204?224.
Establishing conjecture T1.
Louis H. Kauffman, State models and the Jones polynomial, Topology 26 (1987), no. 3, 395-407. 2. Kunio Murasugi, Jones polynomials and classical conjectures in knot theory, Topology 26 (1987), no. 2, 187-194. 3. Morwen B. Thistlethwaite, A spanning tree expansion of the Jones polynomial, Topology 26 (1987), no. 3, 297-309. 12 WILLIAM W. MENASCO
Vaughan Jones, A Polynomial Invariant for Knots via von Neumann Algebras. Bull. Am. Math. Soc. 12, 103-111, 1985. 5. Vaughan Jones, ”Hecke Algebra Representations of Braid Groups and Link Polynomials.” Ann. Math. 126, 335-388, 1987. 5. W.B.R. Lickorish, An introduction to knot theory, Graduate Texts in Math-ematics Series, Springer Verlag, 1997.
Establishing conjectures T2 & T3.
William W. Menasco and M. B Thistlethwaite, The classification of alternating links, Ann. Math. 138 (1993), 113-171. 2. William W. Menasco and M. B Thistlethwaite, The Tait flyping conjecture. Bull. Amer. Math. Soc. (N.S.) 25 (1991), no. 2, 403-412.
What is an alternating knot?
Joshua A. Howie, A characterisation of alternating knot exteriors. Geom. Topol. 21 (2017), no. 4, 2353-2371. 2. Joshua E. Greene, Alternating links and definite surfaces. With an appendix by Andrs Juhsz and Marc Lackenby. Duke Math. J. 166 (2017), no. 11, 2133-2151. |
13740 | https://www.liceomalpighi.it/didattica/mfilipucci/downloads/5AL-5BLOndeSuono_Cutnell_Zanichelli.pdf | La musica è bella ma può essere pericolosa. La prima fila del pubblico a un concerto rock si trova a circa 4 m di distanza dall’altoparlante più vicino, che emette una potenza sonora media di 100 W.
L’ordine di grandezza A quanti decibel è esposto uno spettatore della prima fila? 331 12 C A P I T O L O La risposta a pagina 355 Le onde e il suono 12.1 La natura delle onde Le onde sull’acqua hanno due caratteristiche in comune con tutte le altre onde: • un’onda è una perturbazione che si propaga nello spazio; • un’onda trasporta energia da un posto a un altro.
Nella figura 12.1 l’onda generata dal motoscafo si muove nel lago e disturba il pescatore, perché trasferisce alla barca parte della sua energia. L’onda, però, non trasporta una massa d’acqua ed è quindi molto diversa dalla corrente in un fiume; piuttosto è una perturbazione che si propaga sulla superficie del lago.
Onde trasversali Supponiamo di appoggiare su un tavolo una lunga molla elicoidale. Se muoviamo velocemente un estremo della molla in direzione perpendicolare a essa, vediamo che si propaga lungo la molla una perturbazione che chiamiamo impulso (figura 12.2).
Se muoviamo l’estremo della molla con moto periodico perpendicolare a essa, lun-go la molla si propaga una perturbazione detta onda periodica. In particolare, se questo moto è armonico sulla molla si propaga un’onda armonica (figura 12.3).
Come mostra la parte C della figura, l’onda è formata da una successione di regio-ni della molla che oscillano alternativamente verso l’alto e verso il basso e si pro-pagano verso destra trasferendo la perturbazione a ogni spira della molla. Fissan-do un segno colorato in un punto della molla, si nota che il segno compie oscilla-zioni armoniche attorno alla sua posizione di equilibrio e che i suoi spostamenti sono perpendicolari, o «trasversali», alla direzione in cui si propaga l’onda.
CAPITOLO 12 Le onde e il suono 332 A B C Figura 12.1 L’onda prodotta dal motoscafo si propaga sulla superficie del lago e disturba il pescatore.
Figura 12.3 Quando si fa oscillare rapidamente e con continuità verso l’alto e verso il basso un estremo della molla, si genera un’onda trasversale che si propaga da sinistra verso destra.
Figura 12.2 Muovendo rapidamente la molla in direzione perpendicolare al suo asse, si provoca una deformazione (impulso) che si propaga lungo la molla stessa da sinistra verso destra.
A B C Le onde radio, le onde luminose e le microonde sono esempi di onde trasversali.
Sono onde trasversali anche quelle che si propagano nelle corde di strumenti musi-cali come la chitarra o il violino.
Onde longitudinali Consideriamo nuovamente una molla elicoidale appoggiata su un tavolo. Se muo-viamo avanti e indietro con moto armonico un estremo della molla nella direzione della sua lunghezza (cioè longitudinalmente), come nella figura 12.4, lungo la mol-la si propaga un’onda. Come mostra la parte C della figura, l’onda è formata da una successione di regioni della molla alternativamente compresse e diradate che si pro-pagano lungo la molla. In questo caso, un segno colorato fissato sulla molla si muo-ve avanti e indietro compiendo oscillazioni armoniche nella stessa direzione in cui si propaga l’onda.
ONDA TRASVERSALE Un’onda è trasversale quando la direzione della perturbazione è perpendico-lare alla direzione di propagazione dell’onda.
Le onde sonore, per esempio, sono onde longitudinali.
Alcuni tipi di onde non sono né trasversali né longitudinali. Per esempio, in un’onda che si propaga sulla superficie dell’acqua le particelle non si spostano né in direzione perpendicolare a quella in cui viaggia l’onda né nella stessa direzio-ne. Come mostra la figura 12.5, i loro spostamenti hanno infatti sia una compo-nente perpendicolare sia una componente parallela alla direzione di propagazio-ne dell’onda. In particolare, le particelle d’acqua più vicine alla superficie descri-vono traiettorie quasi circolari.
12.2 Onde periodiche Le onde periodiche sono costituite da cicli, cioè da fenomeni che si ripetono ugua-li a se stessi con un ritmo regolare. Per esempio, un ciclo completo delle onde del-le figure 12.3 e 12.4 è costituito da un’oscillazione armonica di ogni punto della molla intorno alla sua posizione di equilibrio.
ONDA LONGITUDINALE Un’onda è longitudinale quando la direzione della perturbazione è uguale alla direzione di propagazione dell’onda.
CAPITOLO 12 Le onde e il suono 333 Figura 12.4 Quando si fa oscillare rapidamente e con continuità avanti e indietro un estremo della molla, si genera un’onda longitudinale che si propaga da sinistra verso destra.
Figura 12.5 Le onde sull’acqua non sono né trasversali né longitudinali, perché le particelle vicino alla superficie dell’acqua descrivono traiettorie quasi circolari.
B A C Regione di compressione Regione di compressione Regione di diradazione Direzione di propagazione dell’onda Una particella vicina alla superficie dell’acqua descrive una traiettoria circolare quando è investita da un’onda Componente trasversale Componente longitudinale La figura 12.6A mostra la rappresentazione spaziale dell’onda, o forma d’onda, e corrisponde a una «fotografia» dell’onda in un dato istante. Un ciclo è la parte in colore: l’onda è formata da una successione di molti cicli.
L’ampiezza A dell’onda è lo spostamento massimo dalla posizione d’equilibrio di un punto del mezzo in cui si propaga l’onda. L’ampiezza dell’onda è uguale alla distanza tra una cresta (il punto più alto) o un ventre (il punto più basso) e la posi-zione di equilibrio.
La lunghezza d’onda λ (lambda) è la lunghezza di un ciclo ed è uguale alla distanza tra due creste successive, o tra due ventri successivi.
La figura 12.6B è invece una rappresentazione temporale dell’onda e mostra come varia la posizione di un dato punto del mezzo in cui si propaga l’onda al varia-re del tempo. Al passaggio dell’onda, il punto osservato compie oscillazioni armo-niche intorno alla posizione di equilibrio.
Il periodo T è l’intervallo di tempo in cui viene compiuta un’oscillazione com-pleta. In modo equivalente, si può definire il periodo come il tempo impiegato dal-l’onda per percorrere una distanza uguale alla lunghezza d’onda.
Il periodo e la frequenza sono legati dalla relazione: f Come abbiamo già visto nel capitolo 4, il periodo si misura in s e la frequenza in s1, cioè in hertz (Hz). Per esempio, se un’onda impiega un decimo di secondo per compiere un intero ciclo, in ogni secondo compie 10 cicli e quindi la sua frequen-za è: f 10 s1 10 Hz Il periodo, la lunghezza d’onda e la velocità di propagazione dell’onda sono legati da una relazione molto semplice che possiamo ricavare facendo riferimento alla figura 12.7. Immagina di essere fermo a un passaggio a livello e di osservare il pas-saggio di un treno che viaggia a velocità v. Il treno è formato da vagoni identici tra loro,ciascuno dei quali ha una lunghezza λ e impiega un tempo T per passarti davan-ti: la velocità del treno è perciò v λ/T.
In modo analogo, si dimostra che la velocità di propagazione di un’onda di lun-ghezza d’onda λ e periodo T è: v f λ (12.1) La terminologia che abbiamo introdotto e le relazioni fondamentali f 1/T e v f λ valgono sia per le onde trasversali sia per le onde longitudinali.
λ T 1 0,1 s 1 T CAPITOLO 12 Le onde e il suono 334 Figura 12.7 Un treno che viaggia a velocità costante può essere utilizzato come analogia per un’onda che si propaga con la stessa velocità.
Figura 12.6 Le parti evidenziate in colore rappresentano un ciclo d’onda.
L’ampiezza dell’onda è indicata con A.
Lunghezza d’onda = λ λ Posizione verticale sulla molla Posizione verticale di un punto della molla Posizione di equilibrio In un dato istante In un dato punto Distanza A A Tempo Periodo = T A A A B Velocità = v Il tempo impiegato da un vagone per attraversare l’incrocio è il periodo T Lunghezza d’onda = λ λ La velocità di un’onda su una corda La velocità di propagazione di un’onda dipende dalle caratteristiche del mezzo in cui si propaga ().
In una corda sottoposta a una tensione T, la velocità di propagazione di un’on-da trasversale aumenta all’aumentare di T e diminuisce all’aumentare della massa per unità di lunghezza m/L della corda. Si può infatti dimostrare che, per piccole ampiezze, la velocità di propagazione di un’onda trasversale in una corda è: v (12.2) Negli strumenti musicali a corda, come la chitarra, il violino e il pianoforte, le onde trasversali sono generate pizzicando le corde, strofinandole con un archetto o per-cuotendole con un martelletto.
T m/L CAPITOLO 12 Le onde e il suono 335 ESEMPIO 1 Onde trasversali La lunghezza d’onda delle onde radio Le onde radio AM e FM sono onde trasversali costituite da perturbazioni di tipo elettromagnetico che si propagano a 3,00 108 m/s e che tratteremo nel capitolo 24. Una stazione radio trasmette onde radio AM con una frequenza di 1230 kHz e onde radio FM con una frequenza di 91,9 MHz.
Calcola la lunghezza d’onda di ciascuno dei due tipi di onde.
Ragionamento e soluzione Ricordando che 1 kHz 103 Hz e 1 MHz 106 Hz, calcoliamo le lunghezze d’on-da mediante la relazione v f λ: AM λ FM λ Osserva che la lunghezza d’onda di un’onda radio AM è uguale a circa due vol-te e mezzo la lunghezza di un campo da calcio!
3,26 m 3,00 108 m/s 91,9 106 Hz v f 244 m 3,00 108 m/s 1230 103 Hz v f ESEMPIO 2 Velocità di un’onda trasversale Velocità delle onde sulle corde di una chitarra Quando si pizzica la corda di una chitarra, su di essa si propagano onde trasver-sali (figura 12.8). La lunghezza delle corde fra i due estremi fissi è 0,628 m. La massa della corda del mi acuto è 0,208 g, mentre quella del mi grave è 3,32 g. In entrambe la tensione è 226 N.
Calcola le velocità di propagazione delle onde nelle due corde.
Ragionamento e soluzione Poiché la tensione è uguale in entrambe le corde, ci aspettiamo che la velocità dell’onda sia maggiore nella corda che ha densità lineare minore. Infatti: Mi acuto v Mi grave v Nota quanto sono elevate queste velocità: i due valori calcolati equivalgono rispettivamente a 2970 km/h e 745 km/h.
207 m/s 226 N (3,32 103 kg)/(0,628 m) T m/L 826 m/s 226 N (0,208 103 kg)/(0,628 m) T m/L Problem solving Osservazione sulla velocità di un’onda La relazione v f λ vale per qualunque tipo di onda periodica.
() Le onde elettromagnetiche, che verranno trattate nel capitolo 24, possono propagarsi anche nel vuo-to, oltre che in materiali come il vetro e l’acqua.
Figura 12.8 Pizzicando le corde di una chitarra si generano onde trasversali che si propagano nelle corde.
Vibrazione trasversale della corda Fisica quotidiana Onde sulle corde di una chitarra 12.3 La descrizione matematica di un’onda Quando un’onda si propaga in un mezzo sposta le particelle del mezzo dalla loro posizione di equilibrio. Nel caso di un’onda periodica generata dal moto armonico della sorgente, lo spostamento è espresso mediante una funzione goniometrica del tempo, come il seno o il coseno. Ciò era prevedibile, perché nel capitolo 4 abbiamo visto che il moto armonico si descrive mediante funzioni goniometriche e il grafico dello spostamento in funzione del tempo ha l’aspetto mostrato in figura 12.6.
Consideriamo una particella a distanza x dall’origine del sistema di riferimento; lo spostamento y di questa particella in ogni istante t in cui transita un’onda di ampiezza A, frequenza f e lunghezza d’onda λ è Onda che si propaga verso x y A sen 2πft (12.3) Onda che si propaga verso x y A sen 2πft (12.4) Le equazioni precedenti si applicano a onde trasversali o longitudinali tali che y 0 m quando x 0 m e t 0 s. L’angolo 2πft è detto fase dell’on-da e si deve misurare in radianti.
Per comprendere il significato dell’equazione (12.3), consideriamo un’onda tra-sversale su una corda che si propaga verso x, cioè verso x crescenti. Una particella della corda posta nell’origine (x 0 m) si muove di moto armonico con una fase 2πft; quindi il suo spostamento in funzione del tempo è dato da: y A sen (2πft) Una particella posta a una distanza x si muove di moto armonico con una fase: 2πft 2πf t 2πf t La grandezza x/v è il tempo che l’onda impiega per coprire la distanza x. Quindi il moto armonico del punto della corda a distanza x dall’origine è ritardato di un inter-vallo di tempo x/v rispetto al moto armonico nell’origine. In altre parole: in x un punto della corda si muove con il moto armonico che aveva un punto nell’origine x/v secondi prima.
La figura 12.9 mostra lo spostamento in funzione della posizione lungo la cor-da in una sequenza di istanti di tempo separati da 1/4 di periodo T (t 0 s, T/4, 2T/4, 3T/4, T). I grafici sono calcolati sostituendo i corrispondenti valori di t nell’e-quazione (12.3) e ricordando che f 1/T. Questi grafici rappresentano una serie di fotografie istantanee prese mentre l’onda si propaga verso destra. Il quadratino colorato di ogni grafico segna il punto dell’onda che si trovava nell’origine x 0 m all’istante t 0 s. Al passare del tempo, il quadratino si sposta verso destra.
In modo analogo, si può verificare che l’equazione (12.4) rappresenta un’onda che si propaga verso sinistra, cioè nella direzione negativa delle x.
12.4 La natura del suono Onde sonore longitudinali Il suono è costituito da onde longitudinali generate da un oggetto che vibra, la sor-gente sonora, come la corda di una chitarra, le corde vocali umane o la membra-na di un altoparlante. Il suono può essere generato e trasmesso solo in un mezzo materiale, come un gas, un liquido o un solido. Quindi il suono non si propaga nel vuoto.
Per capire come si generano le onde sonore e perché sono onde longitudinali, consideriamo la membrana di un altoparlante. Come mostra la figura 12.10A, quan-do la membrana si muove verso l’esterno, comprime lo strato d’aria davanti a essa, x v x fλ 2πx λ 2πx λ 2πx λ 2πx λ CAPITOLO 12 Le onde e il suono 336 Figura 12.9 I grafici rappresentano la forma di un’onda che si propaga verso destra e sono stati calcolati a partire dall’equazione (12.3) per gli istanti di tempo di volta in volta indicati. Il quadratino colorato in ciascun grafico segna il punto dell’onda che nell’istante t 0 s si trovava nel punto x 0 m. Al passare del tempo, l’onda si propaga verso destra.
y +x t = 0 s t = T 1 – 4 t = T 2 – 4 t = T 3 – 4 t = T provocando un leggero aumento della pressione dell’aria in questa regione. La regione di aria compressa è chiamata compressione e si allontana dall’altoparlante con la velocità del suono. La compressione è analoga alla regione di spire compresse in un’onda longitudinale che si propaga in una molla.
Dopo aver prodotto la compressione, la membrana dell’altoparlante torna indietro, muovendosi verso l’interno, come mostra la figura 12.10B, e producendo nello stra-to d’aria a contatto con essa una regione chiamata rarefazione, in cui la pressione dell’aria è leggermente minore di quella normale. La rarefazione è analoga alla regione di spire diradate di un’onda longitudinale che si propaga in una molla.
Anche la rarefazione si allontana dall’altoparlante alla velocità del suono seguen-do la compressione.
Continuando a vibrare, la membrana genera una successione di zone in cui l’aria è alternativamente compressa e rarefatta: è questa l’onda sonora. La figura 12.11 mostra che le molecole d’aria, come le particelle della molla, oscillano avanti e indie-tro nella stessa direzione in cui si propaga l’onda: quindi l’onda sonora è longitudi-nale. La figura mostra anche che la lunghezza d’onda λ è uguale alla distanza tra i centri di due compressioni successive o di due rarefazioni successive.
La figura 12.12 mostra come si propaga nello spazio l’onda sonora generata dalla membrana dell’altoparlante. Quando le compressioni e le rarefazioni arriva-no all’orecchio, fanno vibrare il timpano con la stessa frequenza della membrana vibrante. Il moto vibratorio del timpano viene interpretato come suono dal cer-vello. È importante ricordare che il suono non è generato dagli spostamenti di masse d’aria come quelli che danno origine ai venti. Le compressioni e le rarefa-zioni generate dall’altoparlante non trasportano le molecole dell’aria, ma le fanno oscillare attorno alla loro posizione d’equilibrio, a cui ritornano quando l’onda è passata.
Frequenza di un’onda sonora Ogni ciclo di un’onda sonora è composto da una compressione e una rarefazione, e la frequenza dell’onda è il numero di cicli che passano in un secondo in uno stes-so punto del mezzo in cui l’onda si propaga.
Per esempio, se la membrana di un altoparlante oscilla avanti e indietro con moto armonico di frequenza 1000 Hz, essa genera in un secondo una successione di 1000 compressioni seguite da 1000 rarefazioni, cioè genera un’onda sonora che ha una frequenza di 1000 Hz.
CAPITOLO 12 Le onde e il suono 337 vmolla vmolla Compressione Pressione dell’aria normale v Rarefazione Compressione Pressione dell’aria normale v A B Figura 12.11 Sia l’onda che si propaga nella molla sia l’onda sonora generata dalla membrana di un altoparlante sono onde longitudinali. I pallini colorati attaccati a una spira della molla e a una molecola d’aria vibrano avanti e indietro in direzione parallela a quella in cui si propaga l’onda.
Figura 12.12 Le regioni di compressione e di rarefazione si propagano nell’aria dall’altoparlante verso l’orecchio dell’ascoltatore, ma le molecole dell’aria non sono trasportate dall’onda. Ogni molecola vibra avanti e indietro intorno alla sua posizione di equilibrio.
Figura 12.10 A. Quando la membrana vibrante di un altoparlante si muove verso l’esterno, produce una compressione dello strato d’aria a contatto con essa.
B. Quando la membrana si muove verso l’interno, produce una rarefazione dello strato d’aria a contatto con essa.
Le regioni di compressione e rarefazione sono simili a quelle di un’onda longitudinale che si propaga in una molla.
Lunghezza d’onda λ = v vmolla Fisica quotidiana La membrana di un altoparlante Vibrazione di una molecola d’aria Un suono si dice puro quando le particelle investite dall’onda sonora oscillano con moto armonico (figura 12.13A): la frequenza di oscillazione delle particelle è la frequenza del suono puro. Quando le particelle oscillano di moto periodico ma non armonico (figura 12.13B), il suono si dice complesso: anche in questo caso, si può individuare una frequenza che caratterizza il suono, detta frequenza fonda-mentale.
Una persona giovane riesce a sentire suoni che hanno frequenze comprese tra 20 Hz e 20 000 Hz (cioè 20 kHz). La capacità di sentire i suoni con frequenza mag-giore diminuisce con l’età: una persona di mezz’età non riesce più a percepire suo-ni con frequenze superiori a 12-14 kHz.
È possibile generare suoni che hanno frequenze minori o maggiori dei limiti di udi-bilità, anche se normalmente questi suoni non vengono percepiti dall’orecchio uma-no. I suoni con frequenza minore di 20 Hz sono chiamati infrasuoni, mentre quelli con frequenza maggiore di 20 kHz sono chiamati ultrasuoni. I rinoceronti (figura 12.14A) si chiamano tra loro emettendo infrasuoni con frequenza di circa 5 Hz, mentre i pipistrelli (figura 12.14B) usano ultrasuoni con frequenze fino a 100 kHz per individuare la posizione degli oggetti e per orientarsi nel volo.
Un suono puro può essere generato con un diapason (figura 12.15), mentre gli stru-menti musicali non emettono suoni puri ma suoni complessi caratterizzati da for-me d’onda molto differenti tra loro (figura 12.16).
Altezza e timbro La frequenza è una caratteristica oggettiva del suono perché può essere misurata con un apposito strumento. Invece il modo in cui la frequenza viene percepita cam-bia da un ascoltatore all’altro. Il nostro cervello, infatti, interpreta le frequenze rile-vate dall’orecchio in termini di una qualità soggettiva detta altezza: un suono con una frequenza fondamentale alta è interpretato come un suono alto o acuto, men-tre un suono con una frequenza fondamentale bassa è interpretato come un suono basso o grave. Per esempio, un ottavino produce suoni acuti, mentre una tuba pro-duce suoni gravi.
CAPITOLO 12 Le onde e il suono 338 Figura 12.14 A. I rinoceronti si chiamano tra loro usando infrasuoni.
B. I pipistrelli usano ultrasuoni per orientarsi nel volo e per individuare la posizione degli oggetti.
Figura 12.15 Il diapason è uno strumento sonoro formato da una forcella d’acciaio a due rami (rebbi) che, percossi, emettono un suono puro. Del diapason riparleremo più diffusamente nel paragrafo 12.9.
Figura 12.13 La rappresentazione temporale di un suono puro (A) e di un suono complesso (B). In questo caso la frequenza fondamentale del suono complesso è la stessa del suono puro.
y Tempo y Tempo A B A B © Ron Garrison, Zoological Society, San Diego © Merlin D. Tuttle/Bat Conservation International/Photo Researchers © Elmtree Images / Alamy Le note della scala musicale corrispondono a ben precise frequenze sonore. Senza entrare nel dettaglio della notazione musicale, ci limitiamo a riportare in figura 12.17 le frequenze dei diversi do nella tastiera di un pianoforte.
Nel complesso, il nostro udito è uno strumento assai raffinato. Infatti, quando ascoltiamo una stessa nota musicale suonata da strumenti diversi, siamo in grado di distinguerli anche se gli strumenti stanno emettendo suoni con la stessa frequenza fondamentale. Questa diversa percezione è legata a una caratteristica del suono, chiamata timbro, che dipende dalla particolare legge periodica con cui oscillano le particelle quando sono investite dall’onda sonora.
L’ampiezza di un’onda sonora La figura 12.18 rappresenta un suono puro che si propaga all’interno di un tubo.
Alcuni manometri disposti lungo il tubo misurano i valori della pressione in vari punti dell’onda. Il grafico che si ottiene riportando i valori della pressione al varia-re della posizione del punto sorgente è una sinusoide. Anche se questo grafico è simile a quello della forma d’onda di un’onda trasversale, ricordiamo che un’onda sonora è un’onda longitudinale e che il grafico posizione-pressione non va confuso con il grafico posizione-spostamento.
CAPITOLO 12 Le onde e il suono 339 Figura 12.16 Ogni strumento musicale ha un proprio timbro a cui corrisponde un tipo particolare di onda periodica. La figura riporta le forme d’onda sonora relative a una stessa nota (sol) emessa da un diapason, da un vibrafono e da un trombone. Le tre onde hanno lo stesso periodo di 2,6 ms, ma forme molto diverse.
Figura 12.17 Frequenze dei diversi do nella tastiera di un pianoforte.
Figura 12.18 Un’onda sonora è costituita da una successione di regioni di compressione e regioni di rarefazione. Il grafico mostra che la pressione dell’aria nelle regioni di compressione è maggiore della pressione atmosferica ed è minore della pressione atmosferica nelle regioni di rarefazione (AP , alta pressione; BP , bassa pressione).
32 64 128 256 512 1024 2048 4096 Frequenze (Hz) BP AP BP AP Pressione dell’aria Alta Posizione Pressione atmosferica normale Rarefazione Bassa Compressione Ampiezza Pressione sonora 0 0 1 2 3 4 6 0 0 1 2 3 4 6 0 0 1 2 3 4 5 6 T 2T Diapason Vibrafono Trombone 5 5 Istante (ms) Il grafico mostra che la pressione è maggiore della pressione atmosferica nelle regio-ni di compressione e minore della pressione atmosferica nelle regioni di rarefazione.
Il grafico mostra anche l’ampiezza della pressione, che è il valore massimo della dif-ferenza tra la pressione in una regione di compressione e il valore della pressione atmosferica normale.Anche in questo caso osserviamo che l’ampiezza della pressio-ne non va confusa con l’ampiezza dell’onda.
Le variazioni di pressione in un’onda sonora sono in genere molto piccole.
Per esempio, in una normale conversazione tra due persone l’ampiezza della pres-sione è di circa 3 102 Pa, ben minore della pressione atmosferica normale che è pari a 1,013 105 Pa.
La caratteristica del suono che dipende dall’ampiezza della pressione è la sua intensità: tanto maggiore è l’ampiezza della pressione, tanto più forte è il suono.
L’intensità di un suono è una caratteristica sia oggettiva sia soggettiva. Infatti l’am-piezza della pressione può essere misurata con opportuni strumenti, ma uno stesso suono può essere percepito come più forte o più debole da persone diverse a secon-da della sensibilità del loro apparato uditivo.
La velocità del suono Come mostra la tabella 12.1, il suono si propaga a velocità molto diverse nei gas, nei liquidi e nei solidi. A temperatura ambiente la velocità del suono nell’aria è di 343 m/s (pari a circa 123018 km/h) ed è notevolmente maggiore nei liquidi e nei solidi. Per esempio, la velocità del suono nell’acqua è circa quattro volte maggiore di quella nell’aria e la sua velocità nell’acciaio è più di diciassette volte maggiore di quella nell’aria. In generale, il suono viaggia più lentamente nei gas, più velocemente nei liquidi e ancora più velocemente nei solidi.
12.5 L’intensità del suono Le onde sonore trasportano energia che può essere usata per compiere lavoro, per esempio per far vibrare il timpano del nostro orecchio. La quantità di energia tra-sportata in un secondo da un’onda è chiamata potenza dell’onda e nel Sistema Inter-nazionale si misura in joule al secondo (J/s), cioè in watt (W).
Intensità di un suono Quando si allontana dalla sorgente che l’ha emessa, come l’altoparlante della figu-ra 12.19, un’onda sonora si propaga attraversando superfici di area sempre mag-giore. Attraverso le superfici indicate in figura con 1 e 2 passa la stessa potenza, ma l’intensità del suono è minore nella superficie più lontana.
INTENSITÀ DI UN SUONO L’intensità di un suono I è definita come rapporto tra la potenza sonora media – P che attraversa perpendicolarmente una data superficie e l’area A della super-ficie: I — – P A (12.5) Unità di misura: watt al metro quadrato (W/m2).
1 2 CAPITOLO 12 Le onde e il suono 340 Figura 12.19 La potenza trasportata da un’onda sonora si propaga con l’onda dopo che è stata emessa dalla sorgente, che in questa figura è l’altoparlante.
Perciò l’onda attraversa perpendicolarmente prima la superficie 1 e poi la superficie 2 che ha un’area maggiore.
Tabella 12.1 Velocità del suono in alcuni sostanze gassose, liquide e solide Sostanza Velocità (m/s) Gas Aria (a 0 °C) 331 Aria (a 20 °C) 343 Biossido di carbonio (a 0 °C) 259 Elio (a 0 °C) 965 Ossigeno (a 0 °C) 361 Liquidi Acqua dolce (a 20 °C) 1482 Acqua di mare (a 20 °C) 1522 Alcol etilico (a 20 °C) 1162 Cloroformio (a 20 °C) 1004 Mercurio (a 20 °C) 1450 Solidi Acciaio 5960 Piombo 5010 Rame 1960 Vetro (Pyrex) 5640 Se una sorgente emette onde sonore in modo uniforme in tutte le direzioni, l’in-tensità sonora è legata alla distanza da una relazione molto semplice. La figura 12.20 mostra una sorgente di questo tipo al centro di una sfera immaginaria (di cui per chiarezza è rappresentata solo una metà) di raggio r. Poiché l’intera superficie sferica (di area A 4πr2) è attraversata dalla stessa potenza P, l’intensità del suo-no I in un punto a distanza r dalla sorgente è: Onda sferica uniforme I (12.6) Dall’equazione precedente si può vedere che l’intensità di un’onda sferica unifor-me è inversamente proporzionale al quadrato della distanza dalla sorgente che l’ha emessa. Per esempio, se la distanza raddoppia l’intensità del suono diventa un quar-to: (1/2)2 1/4.
P 4πr2 CAPITOLO 12 Le onde e il suono 341 Figura 12.20 La sorgente sonora al centro della sfera (di cui è rappresentata per chiarezza solo una metà) emette onde sonore uniformemente in tutte le direzioni.
Figura 12.21 Se il suono prodotto dallo scoppio di un fuoco d’artificio si propaga uniformemente in tutte le direzioni, la sua intensità in un punto a distanza r dalla sorgente è I P/4r2, dove P è la potenza sonora generata dallo scoppio.
ESEMPIO 3 Intensità di un suono Suoni da un altoparlante L’altoparlante della figura 12.19 genera una potenza sonora di 12 105 W che attraversa perpendicolarmente le superfici indicate con 1 e 2, di area rispettiva-mente A1 4,0 m2 e A2 12 m2.
Calcola l’intensità del suono in ciascuna delle due superfici e spiega perché l’os-servatore 2 sente un suono più debole di quello che sente l’osservatore 1.
Ragionamento e soluzione Le superfici 1 e 2 sono attraversate dalla stessa potenza, ma l’area della superfi-cie 2 è maggiore di quella della superficie 1, quindi l’intensità del suono nella superficie 2 è minore di quella nella superficie 1. Dall’equazione (12.5): Superficie 1 I1 Superficie 2 I2 L’orecchio di un osservatore, che ha sempre la stessa superficie, è attraversato da una potenza minore dove l’intensità del suono (cioè la potenza sonora per unità di superficie) è minore. Quindi l’osservatore 2 sente un suono più debole.
1,0 10–5 W/m2 12 105 W 12 m2 – P A2 3,0 10–5 W/m2 12 105 W 4,0 m2 – P A1 ESEMPIO 4 Onde sferiche Fuochi d’artificio La figura 12.21 mostra l’esplosione di un fuoco d’artificio. Supponi che il suono generato dallo scoppio si propaghi uniformemente in tutte le direzioni e che il suo-no riflesso dal suolo sia trascurabile. Quando il suono arriva all’osservatore 2, che si trova a una distanza r2 640 m dalla sorgente,la sua intensità è I2 0,10W/m2.
r2 r1 2 1 Problem solving Intensità sonora e potenza sonora L’intensità sonora I e la potenza sonora P sono due grandezze diverse e non vanno confuse. Sono comunque due grandezze collegate, perché la potenza sonora è uguale all’intensità sonora per unità di superficie.
Sorgente sonora al centro della sfera CAPITOLO 12 Le onde e il suono 342 Livello di intensità sonora L’intensità minima I0 di un suono puro con una frequenza di 1000 Hz che può essere percepita da un orecchio umano è I0 1 1012 W/m2. Questo valore del-l’intensità è chiamato soglia minima udibile. All’altro estremo si trova la soglia massima sopportabile (o soglia del dolore), che è circa 10 W/m2, ma già un’espo-sizione continua a suoni di intensità superiore a 1 W/m2 provoca dolori e danni permanenti all’apparato uditivo. È comunque notevole l’ampiezza dell’intervallo di intensità a cui l’orecchio umano è sensibile: l’intensità massima che l’orecchio umano può tollerare senza dolore è mille miliardi di volte l’intensità minima per-cepibile.
La nostra percezione del volume di un suono dipende dall’intensità sonora secondo una relazione che non è lineare ma logaritmica. Per questa ragione si sce-glie di misurare l’intensità con cui percepiamo un suono di intensità I mediante il livello di intensità sonora β, così definito: β 10 log10 dove I0 1 1012 W/m2 è la minima intensità sonora udibile. Nel Sistema Interna-zionale il livello di intensità sonora è espresso in decibel (dB) (tabella 12.2). Espe-rimenti di fisiologia dimostrano che il volume percepito sembra raddoppiare quan-do il livello di intensità sonora aumenta di 10 dB.
Tabella 12.2 Intensità sonore e livelli di intensità sonora per alcune tipologie di suono Intensità Livello di intensità sonora I (W/m2) sonora β (dB) Soglia di udibilità 1,0 1012 0 Fruscio di foglie 1,0 1011 10 Bisbiglio 1,0 1010 20 Conversazione normale (a 1 metro) 3,2 106 65 Interno di un’automobile nel traffico 1,0 104 80 Rumore di un’automobile senza marmitta 1,0 102 100 Concerto rock dal vivo 1,0 102 120 Soglia del dolore 10 102000 130 I I0 Qual è l’intensità del suono che arriva all’osservatore 1 che si trova a una distanza r1 160 m dalla sorgente?
Ragionamento e soluzione Calcoliamo il rapporto tra le intensità sonore percepite dai due ascoltatori mediante l’equazione (12.6): 16 Di conseguenza: I1 16 I2 16 (0,10 W/m2) 1,6 W/m2 (640 m)2 (160 m)2 r2 2 r2 1 P 4πr2 1 P 4πr2 2 I1 I2 Problem solving Osservazione sulle onde sferiche L’equazione (12.6) può essere usata solo quando l’onda sonora si propaga uniformemente in tutte le direzioni e non ci sono ostacoli che possano rifletterla.
12.6 L’effetto Doppler Quando un camion dei pompieri si avvicina a noi il suono della sua sirena è più acu-to di quando il camion è fermo o si allontana. Qualcosa di simile accade quando ci avviciniamo a una sorgente sonora ferma o ci allontaniamo da essa. Questi feno-meni furono spiegati nel 1842 dal fisico austriaco Christian Doppler (1803-1853) e sono chiamati collettivamente «effetto Doppler».
L’effetto Doppler è la variazione di frequenza del suono rilevato dal ricevitore perché la sorgente sonora e il ricevitore hanno velocità diverse rispetto al mezzo in cui il suono si propaga.
Sorgente in movimento e ricevitore fermo Per capire come si verifica l’effetto Doppler nel caso in cui la sorgente sonora è in movimento e il ricevitore è fermo, consideriamo prima il suono emesso dalla sire-na del camion dei pompieri fermo, rappresentato nella figura 12.22A. Supponiamo che, oltre al camion, anche l’aria sia ferma rispetto alla Terra.
Ciascuno degli archi disegnati in blu nella figura rappresenta una regione di com-pressione dell’onda sonora. Poiché le compressioni e le rarefazioni sono simmetri-che rispetto alla sorgente, sia il ricevitore davanti al camion sia il ricevitore dietro il camion ricevono lo stesso numero di compressioni al secondo, quindi i suoni che essi percepiscono hanno la stessa frequenza.
Quando il camion comincia a muoversi, la situazione cambia nel modo rappre-sentato nella parte B della figura. Le compressioni davanti al camion risultano ora più vicine tra loro e ciò provoca una diminuzione della lunghezza d’onda. Questo «addensamento» è dovuto al fatto che, prima di emettere un’altra compressione, il camion «guadagna terreno» rispetto alla compressione precedente.
Poiché le compressioni davanti al camion sono più vicine tra loro, l’osservatore davanti al camion riceve un numero di compressioni al secondo maggiore di quel-lo che riceveva quando il camion era fermo. Questo significa che il suono percepi-to dall’osservatore ha una frequenza maggiore, cioè è più acuto, del suono percepi-to quando il camion era fermo.
Invece le compressioni dietro il camion sono più lontane tra loro quando il camion è in moto rispetto a quando il camion è fermo, e ciò provoca un aumento della lunghezza d’onda. L’allontanamento delle compressioni dietro il camion è dovuto al fatto che il camion «perde terreno» rispetto alle compressioni emesse in precedenza. L’osservatore dietro il camion riceve perciò un numero di compressio-ni al secondo minore di quello che riceveva quando il camion era fermo. Questo significa che il suono percepito dall’osservatore ha una frequenza minore, cioè è più grave, rispetto al suono percepito quando il camion era fermo.
Se la sirena del camion fermo della figura 12.22A emette una compressione nel-l’istante t 0 s, emetterà la compressione successiva nell’istante T, dove T è il perio-do dell’onda. Come mostra la figura 12.23A, la distanza tra queste due compres-sioni è la lunghezza d’onda λ dell’onda sonora emessa dalla sorgente ferma. Quan-do il camion si muove con velocità vs (dove «s» sta per «sorgente» del suono) ver-so un osservatore fermo, la sirena emette compressioni negli istanti t 0 s e t T.
CAPITOLO 12 Le onde e il suono 343 Figura 12.22 A. Quando il camion è fermo, la frequenza del suono emesso dalla sirena è uguale sia davanti al camion sia dietro di esso.
B. Quando il camion è in movimento, la lunghezza d’onda del suono davanti al camion diminuisce mentre quella del suono dietro il camion aumenta.
Figura 12.23 A. Quando il camion è fermo, la distanza tra due regioni di compressione successive è uguale alla lunghezza d’onda λ del suono emesso dalla sirena.
B. Quando il camion si muove con velocità vs, la lunghezza d’onda λ′ del suono davanti al camion è minore di quella del suono emesso dalla sirena.
A Regioni di compressione Lunghezza d’onda Camion fermo B Lunghezza d’onda maggiore Lunghezza d’onda minore Camion in movimento A B Camion in movimento s sT Osservatore fermo Camion fermo Osservatore fermo λ λ′ Tuttavia, prima di emettere la seconda compressione, il camion si è avvicinato al ricevitore spostandosi in avanti di vsT, come si può vedere nella figura 12.23B.
Di conseguenza, la distanza fra due compressioni successive non è più uguale alla lunghezza d’onda λ del suono emesso dalla sorgente ferma, ma è λ′ (minore di) λ′ e dato da: λ′ λ vsT, Indicando con fr la frequenza percepita dal ricevitore (dove «r» sta per «ricevito-re»), dall’equazione (12.1) si ricava che fr è uguale al rapporto tra la velocità del suono v e la lunghezza d’onda λ′: fr La lunghezza d’onda del suono emesso dalla sirena ferma è λ v/fs, dove fs è la fre-quenza dell’onda emessa dalla sorgente, mentre il suo periodo è T 1/fs. Sosti-tuendo questi valori di λ e di T nell’equazione precedente e risolvendo l’equazio-ne rispetto a fr, si ottiene: Sorgente in moto verso un osservatore fermo fr fs (12.7) Poiché il termine 1 vs/v che compare al denominatore della frazione nell’e-quazione (12.7) è minore di uno, la frequenza fr del suono percepito dal ricevito-re è maggiore della frequenza fs del suono emesso dalla sorgente. La differenza fr fs tra queste due frequenze è chiamata spostamento Doppler e il suo valore dipende dal valore del rapporto tra la velocità della sorgente vs e la velocità del suono v.
Quando la sirena, invece di avvicinarsi al ricevitore, si allontana da esso, la lun-ghezza d’onda tra due regioni di compressione successive λ′ è maggiore di λ e il suo valore è dato da: λ′ λ vsT Applicando lo stesso procedimento che abbiamo usato per ricavare l’equazione (12.7), possiamo trovare il valore fr della frequenza percepita dal ricevitore in que-sto caso: Sorgente che si allontana da un osservatore fermo fr fs (12.8) Poiché il termine 1 vs/v che compare al denominatore della frazione nell’equa-zione (12.8) è maggiore di uno, la frequenza fr del suono percepito dal ricevitore è minore della frequenza fs del suono emesso dalla sorgente.
1 1 — vs v 1 1 — vs v v λ vsT v λ′ CAPITOLO 12 Le onde e il suono 344 ESEMPIO 5 Effetto Doppler Il fischio di un treno che passa Mentre un treno ad alta velocità sta viaggiando a una velocità di 44,7 m/s (161 km/h) il macchinista attiva l’avvisatore acustico che emette un fischio di fre-quenza di 415 Hz. La velocità del suono è 343 m/s. Quali sono le frequenze e le lunghezze d’onda dei suoni percepiti da una persona ferma a un passaggio a livel-lo quando: il treno si sta avvicinando?
il treno si sta allontanando?
Osservatore in movimento e sorgente ferma La figura 12.24 mostra come si verifica l’effetto Doppler quando la sorgente sonora è ferma e il ricevitore è in movimento,sempre nell’ipotesi che anche l’aria sia ferma.
Quando il ricevitore si muove con velocità di modulo vr verso la sorgente fer-ma, percorre una distanza vrt in un tempo t. Durante questo tempo incontra tutte le regioni di compressione che avrebbe incontrato se fosse stato fermo più un cer-to numero di zone di compressione che è dato dal rapporto v rt/λ tra la distanza vrt che ha percorso e la distanza λ tra due regioni di compressione successive. Il nume-ro di regioni di compressione in più che incontra in un secondo è perciò vr/λ. Poi-ché un osservatore fermo sentirebbe un suono con la frequenza fs emessa dalla sor-gente, il ricevitore in moto sente un suono con una frequenza maggiore fr data da: fr fs fs 1 Ricordando che v fsλ, si ottiene: Osservatore che si avvicina a una sorgente ferma fr fs1 (12.9) Un osservatore che si allontana da una sorgente ferma si muove nella stessa dire-zione dell’onda sonora emessa dalla sorgente e quindi incontra un numero minore di regioni di compressione rispetto a quelle che incontrerebbe un osservatore fer-mo. In questo caso il ricevitore sente quindi un suono con una frequenza minore fr data da: Osservatore che si allontana da una sorgente ferma fr fs 1 (12.10) Il meccanismo fisico che produce l’effetto Doppler quando si muove il ricevitore è diverso da quello che lo produce quando si muove la sorgente. Quando si muove la sorgente e il ricevitore è fermo, la lunghezza d’onda λ rappresentata nella figura 12.23B cambia e di conseguenza cambia la frequenza fr del suono percepito dal ricevitore. Quando invece si muove il ricevitore e la sorgente è ferma, la lunghezza d’onda λ rappresentata nella figura 12.24 non cambia. Quello che cambia è il nume-ro di regioni di compressione al secondo incontrate da un osservatore in moto rispetto a quelle incontrate da un osservatore fermo e ciò fa sì che la frequenza fr del suono percepito dal ricevitore in moto sia diversa.
vr v vr v vr fsλ vr λ CAPITOLO 12 Le onde e il suono 345 Ragionamento e soluzione Quando il treno si avvicina, la frequenza del suono percepito dalla persona ferma è: fr fs (415 Hz) La lunghezza d’onda di questo suono è: λ′ Quando il treno si allontana, la frequenza del suono percepito dalla persona ferma è: fr fs (415 Hz) La lunghezza d’onda di questo suono è: λ′ 0,935 m 343 m/s 367 Hz v fr 367 Hz 1 44,7 m/s 1 343 m/s 1 1 — vs v 0,719 m 343 m/s 477 Hz v fr 477 Hz 1 44,7 m/s 1 343 m/s 1 1 — vs v Sorgente ferma Osservatore in movimento rt r λ Figura 12.24 Un osservatore che si muove con velocità vr verso una sorgente sonora ferma incontra un numero maggiore di regioni di compressione al secondo rispetto a un osservatore fermo.
Caso generale Può succedere che sia il ricevitore sia la sorgente si muovano rispetto al mezzo in cui si propaga il suono. Se il mezzo è fermo, la frequenza fr del suono percepito dal ricevitore è data da una combinazione delle equazioni (12.7) e (12.10): Osservatore e sorgente in moto relativo fr fs (12.11) Nel numeratore dell’equazione precedente si deve usare il segno più quando il rice-vitore si muove verso la sorgente e il segno meno quando il ricevitore si allontana dalla sorgente. Nel denominatore si deve usare il segno meno quando la sorgente si muove verso il ricevitore e il segno più quando la sorgente si allontana dal rice-vitore. I simboli vr, vs e v indicano misure senza segno perché i segni più e meno che compaiono nell’equazione tengono già conto dei versi della direzione del moto del ricevitore e della direzione di propagazione dell’onda.
12.7 Il principio di sovrapposizione Succede spesso che in uno stesso punto giungano contemporaneamente due o più onde sonore, come capita per esempio quando più persone parlano tra loro o quan-do una musica arriva dalle casse acustiche di un impianto stereo. Per capire che cosa succede in queste situazioni osserviamo le figure 12.25 e 12.26, che mostrano due impulsi trasversali di ampiezza uguale che viaggiano uno verso l’altro su una molla.
Nella figura 12.25 entrambi gli impulsi sono perturbazioni costituite da sposta-menti verso l’alto delle spire della molla, mentre nella figura 12.26 un impulso è verso l’alto e l’altro verso il basso. Le parti A delle due figure mostrano che cosa succede quando i due impulsi incominciano a sovrapporsi. I due impulsi si combi-nano senza disturbarsi e la molla assume una forma che è la somma delle due for-me che avrebbe avuto al passaggio dei due impulsi separati. Perciò, quando i due impulsi della figura 12.25A si sovrappongono completamente, si verifica la situa-zione mostrata nella parte B della figura, cioè l’ampiezza dell’impulso è il doppio di quella dei due impulsi separati. Analogamente, quando i due impulsi della figu-ra 12.26A si sovrappongono completamente, si verifica la situazione mostrata nel-la parte B della figura, cioè la somma dei due impulsi è temporaneamente nulla e la molla si appiattisce. In entrambi i casi, dopo la sovrapposizione i due impulsi con-tinuano a viaggiare separatamente e la molla si comporta come si comporterebbe se fosse attraversata da due impulsi separati.
Il fatto che due impulsi singoli si sommino senza disturbarsi dando luogo a un impul-so risultante esemplifica una legge più generale chiamata principio di sovrapposizione.
Questo principio vale fra l’altro per le onde sonore, le onde sull’acqua e le onde elettromagnetiche come la luce.
12.8 Interferenza e diffrazione di onde sonore Interferenza Supponiamo che due onde sonore provenienti da due altoparlanti si sovrapponga-no in un punto al centro della zona di ascolto, come nella figura 12.27, e che esse PRINCIPIO DI SOVRAPPOSIZIONE Quando due o più onde sono presenti contemporaneamente in uno stesso pun-to, la perturbazione in quel punto è la somma delle perturbazioni prodotte dal-le singole onde.
1 — vr v 1 — vs v CAPITOLO 12 Le onde e il suono 346 Figura 12.25 Due impulsi trasversali rivolti verso l’alto che si attraversano.
Inizio della sovrapposizione A Sovrapposizione totale B Gli impulsi si allontanano C Figura 12.26 Due impulsi trasversali, uno rivolto verso l’alto e uno rivolto verso il basso, che si attraversano.
Sovrapposizione totale B Gli impulsi si allontanano C Inizio della sovrapposizione A Se i due altoparlanti sono alla stessa distanza dal punto in cui le onde si sovrap-pongono, si ha sempre la sovrapposizione delle compressioni (C) di un’onda con le compressioni dell’altra e la sovrapposizione delle rarefazioni (R) di un’onda con le rarefazioni dell’altra.Per il principio di sovrapposizione la forma dell’onda che risul-ta dalla combinazione delle due onde è uguale alla somma delle forme delle due onde. Di conseguenza, l’ampiezza delle variazioni di pressione nel punto di sovrap-posizione è il doppio dell’ampiezza A delle due onde singole: in questo punto il suo-no è più forte di quello che proviene da uno solo degli altoparlanti. Quando due onde si combinano in modo che una compressione si sovrappone sempre a una com-pressione e una rarefazione si sovrappone sempre a una rarefazione, si dice che sono in fase (o in concordanza di fase) e che danno luogo a interferenza costruttiva.
Supponiamo di allontanare dal punto di sovrapposizione l’altoparlante di sini-stra di mezza lunghezza d’onda (figura 12.28). Nel punto di sovrapposizione una compressione proveniente dall’altoparlante di sinistra incontra una rarefazione pro-veniente dall’altoparlante di destra e, analogamente, una rarefazione proveniente da sinistra incontra una compressione proveniente da destra. Per il principio di sovrapposizione, l’ampiezza dell’onda che risulta dalla combinazione di queste due onde è nulla: le rarefazioni prodotte da un’onda compensano esattamente le com-pressioni prodotte dall’altra (). Il risultato è che la pressione dell’aria rimane costante: in quel punto non si sente alcun suono. Quando due onde si combinano in modo che una compressione si sovrappone sempre a una rarefazione, si dice che sono in opposizione di fase e che danno luogo a interferenza distruttiva.
CAPITOLO 12 Le onde e il suono 347 Figura 12.27 Per effetto dell’interferenza costruttiva tra due onde sonore con la stessa ampiezza A, in un punto posto a distanze uguali dalle due casse acustiche che emettono onde in fase (C, compressioni; R, rarefazioni) si sente un suono più intenso (di ampiezza 2A).
Figura 12.28 Le onde sonore emesse dagli altoparlanti sono in fase. Tuttavia, il fatto che l’altoparlante di sinistra sia a una distanza dal punto di sovrapposizione maggiore di mezza lunghezza d’onda rispetto alla distanza dell’altoparlante di destra produce un fenomeno di interferenza distruttiva tra le due onde (C, compressioni; R, rarefazioni).
Pertanto, nel punto di sovrapposizione non si sente alcun suono.
3 m 3 m C R C R Ricevitore Interferenza costruttiva A Pressione Tempo Pressione Tempo A 2A + = Interferenza distruttiva Tempo Pressione A Pressione Tempo A + = 3 m 3+ m C C R C R Ricevitore 1 – 2 () Quando si allontana la cassa acustica di sinistra, l’intensità del suono emesso da questa cassa, e quin-di anche l’ampiezza della pressione nel punto di sovrapposizione, diminuiscono leggermente. Per fare in modo che l’ampiezza della pressione delle due onde che si sovrappongono sia uguale, in questo capito-lo supporremo che la potenza sonora trasmessa dal ricevitore alla cassa acustica di sinistra sia legger-mente superiore a quella trasmessa alla cassa acustica di destra.
abbiano la stessa ampiezza e la stessa frequenza. Per semplicità supponiamo inoltre che la loro lunghezza d’onda sia λ 1 m e che le membrane degli altoparlanti vibri-no in fase, cioè che esse comincino a muoversi verso l’esterno nello stesso istante e tornino verso l’interno nello stesso istante.
L’interferenza distruttiva è alla base di una tecnica molto utile per ridurre l’inten-sità di un suono o di un rumore indesiderato. Per esempio, la figura 12.29 mostra una coppia di cuffie che limitano i rumori esterni. Dentro le cuffie sono inseriti dei piccoli microfoni che ricevono i suoni (rumori) provenienti dall’esterno. I suoni rice-vuti vengono trasformati in segnali che li riproducono esattamente in opposizione di fase. Questi suoni sono inviati agli altoparlanti delle cuffie e producono un feno-meno di interferenza distruttiva con il suono originale: in questo modo il rumore arriva alle orecchie molto attenuato.
Condizioni di interferenza In generale, non conta il cammino percorso da ciascuna onda, quanto la differenza tra i cammini percorsi per arrivare al punto di sovrapposizione.
L’interferenza costruttiva avviene in tutti i punti in cui si sovrappongono due com-pressioni o due rarefazioni (quattro di questi punti sono indicati con pallini rossi nella figura 12.30). Un osservatore che si trovasse in uno qualunque di questi pun-ti sentirebbe un suono più intenso di quello emesso da ciascuna delle due sorgenti separatamente. Viceversa, l’interferenza distruttiva avviene in tutti i punti in cui si sovrappongono una compressione e una rarefazione (due di questi punti sono indi-cati in figura con pallini verdi). Un osservatore che si trovasse in uno qualunque di questi punti non sentirebbe alcun suono.
Nei punti in cui non si ha né interferenza costruttiva né interferenza distruttiva le due onde si combinano dando luogo a suoni di intensità un po’ maggiore oppu-re un po’ minore del suono emesso da ciascun altoparlante, a seconda delle distan-ze del punto dalle due casse. Pertanto è possibile che un osservatore senta suoni di intensità molto variabile attraversando la regione di sovrapposizione delle onde.
Ciascuna delle due onde sonore della figura 12.30 trasporta energia e l’energia complessiva trasportata nella zona in cui le onde si sovrappongono è la somma del-le energie trasportate da ciascuna onda. Una delle conseguenze più interessanti del CONDIZIONI DI INTERFERENZA DISTRUTTIVA Quando due sorgenti sonore sono in fase, una differenza di cammino rispetto al punto di sovrapposizione pari a mezza lunghezza d’onda o a un numero inte-ro di lunghezze d’onda più mezza lunghezza d’onda — 1 2 , 1 — 1 2 , 2 — 1 2 , 3 — 1 2 ,… produce un fenomeno di interferenza distruttiva.
CONDIZIONI DI INTERFERENZA COSTRUTTIVA Quando due sorgenti sonore sono in fase, una differenza di cammino rispetto al punto di sovrapposizione pari a zero oppure a un numero intero (1, 2, 3,…) di lunghezze d’onda produce un fenomeno di interferenza costruttiva.
CAPITOLO 12 Le onde e il suono 348 Figura 12.30 Due onde sonore si sovrappongono nella zona ombreggiata. Gli archi a tratto pieno indicano le parti centrali delle regioni di compressione (C) mentre quelli tratteggiati indicano le parti centrali delle regioni di rarefazione (R). In ciascuno dei punti indicati con un pallino rosso (•) si ha interferenza costruttiva, mentre in ciascuno dei punti indicati con un pallino verde (•) si ha interferenza distruttiva.
Figura 12.29 Le cuffie antirumore, dette anche cuffie attive, utilizzano il fenomeno dell’interferenza distruttiva.
Fisica quotidiana Le cuffie antirumore Rumore Rumore Altoparlante Microfono Circuito elettronico Rumore in opposizione di fase Livello di rumore attenuato 1 lunghezza d’onda R C R C R C fenomeno dell’interferenza è che l’energia trasportata dalle onde si conserva, ma viene ridistribuita in modo che ci sono zone in cui il suono è più intenso e zone in cui non c’è più alcun suono.
Il fenomeno dell’interferenza riguarda tutti i tipi di onde, non solo le onde sono-re. Parleremo di nuovo dell’interferenza nel capitolo 15, quando tratteremo l’inter-ferenza delle onde luminose.
Diffrazione Quando incontra un ostacolo oppure i bordi di una fenditura, un’onda devia dalla sua direzione di propagazione e prosegue al di là dell’ostacolo o della fenditura. Per esempio, nel superare una porta aperta un’onda sonora si incurva nel modo rap-presentato in figura 12.32A. Se ciò non avvenisse, fuori dalla stanza il suono si sen-tirebbe solo nei punti situati direttamente davanti alla porta, come indica la parte B della figura (supponiamo che l’onda sonora non si propaghi attraverso le pareti della stanza).
CAPITOLO 12 Le onde e il suono 349 ESEMPIO 6 Condizioni di interferenza Che cosa sente un ascoltatore?
I due altoparlanti di figura 12.31 emettono in fase due suoni identici di frequenza 214 Hz e velocità 343 m/s.
L’ascoltatore sente un suono intenso o non sente alcun suono?
Ragionamento e soluzione Per stabilire quale tipo di interferenza si verifica in C, calcoliamo la differen-za AC BC tra i cammini percorsi dalle due onde per arrivare nel punto C.
Poiché ABC è un triangolo rettangolo di cui AC è l’ipotenusa, la distanza AC tra l’altoparlante A e il ricevitore si può calcolare applicando il teorema di Pitagora: AC (3,20 m)2 (2,40 m)2 4,00 m La distanza BC è indicata dal problema come 2,40 m. Perciò la differenza di cam-mino tra le due onde è: AC BC 4,00 m 2,40 m 1,60 m La lunghezza d’onda del suono emesso dai due altoparlanti è: λ 1,60 m Poiché la differenza di cammino è pari a una lunghezza d’onda, nel punto C si ha interferenza costruttiva e l’ascoltatore sente un suono intenso.
343 m/s 214 Hz v f 3,20 m 2,40 m C 90° A B Problem solving Osservazione sulle condizioni di interferenza Per decidere se le onde sonore emesse da due sorgenti interferiscono costruttivamente o distruttivamente in un punto, occorre calcolare la differenza tra i cammini percorsi dalle due onde per arrivare in quel punto e confrontare il risultato ottenuto con la lunghezza d’onda del suono.
Figura 12.31 Sovrapposizione di due onde sonore con una frequenza di 214 Hz.
Con diffrazione Senza diffrazione A B Figura 12.32 A. La curvatura di un’onda sonora attorno ai bordi di una porta è un esempio di diffrazione. La sorgente sonora dentro la stanza non è rappresentata. B. Se non ci fosse il fenomeno della diffrazione, l’onda sonora non si incurverebbe attraversando la porta.
La deviazione dalla direzione di propagazione di un’onda attorno a un ostacolo o ai bordi di una fenditura è chiamato diffrazione. Questo fenomeno riguarda tutti i tipi di onde. Studieremo in maggior dettaglio la diffrazione nel capitolo 15.
Per il momento ci limitiamo a osservare che la diffrazione è provocata dall’in-terferenza delle onde e il suo effetto è quello di far arrivare l’energia trasportata dall’onda in regioni che altrimenti non sarebbero accessibili. Naturalmente l’ener-gia complessiva si conserva: la diffrazione provoca solo una sua ridistribuzione e non si ha né creazione né distruzione di energia.
12.9 Battimenti Il diapason è uno strumento costituito da una forcella d’acciaio che emette un suo-no puro, con una frequenza di 440 Hz, quando viene percosso con un colpo secco.
La figura 12.33 mostra le onde sonore emesse da due diapason identici posti uno accanto all’altro.A uno dei due diapason è però stato attaccato un pezzetto di masti-ce: l’aumento della massa fa diminuire la frequenza del suono emesso da questo diapason a 438 Hz. Quando entrambi i diapason emettono il suono, l’intensità del suono totale aumenta e diminuisce periodicamente: debole, poi forte, poi di nuovo debole e di nuovo forte e così via. Queste variazioni periodiche dell’intensità del suono sono chiamate battimenti e sono dovute all’interferenza di due onde sonore con frequenze leggermente diverse.
Per ragioni di chiarezza nella figura 12.33 le regioni di compressione e di rarefa-zione delle due onde sonore sono rappresentate separatamente, ma in realtà esse si diffondono nello spazio e si sovrappongono e, come previsto dal principio di sovrap-posizione, il suono percepito da un osservatore è la somma dei due suoni. Il suono percepito è forte quando all’orecchio arriva una regione di interferenza costruttiva e debole quando arriva una regione di interferenza distruttiva. Il numero di volte in cui l’intensità del suono passa da forte a debole in un secondo è chiamato fre-quenza dei battimenti.
La frequenza dei battimenti Consideriamo due onde sonore che giungono in un punto fissato, che per sempli-cità scegliamo nell’origine del sistema di riferimento in modo che sia x 0 m. Le due onde sono rappresentate da equazioni analoghe all’equazione (12.3), dove lo spostamento è espresso mediante la funzione coseno: y1 A cos (2πf1t) y2 A cos (2πf2t) CAPITOLO 12 Le onde e il suono 350 Distruttiva Piccolo pezzo di mastice 440 Hz 438 Hz Distruttiva Costruttiva Costruttiva Figura 12.33 Quando i due diapason con frequenze di vibrazione leggermente diverse (440 Hz e 438 Hz) vengono messi in vibrazione contemporaneamente si verifica il fenomeno dei battimenti. Le onde sonore emesse dai due diapason non sono rappresentate in scala.
Per determinare l’ampiezza dell’onda risultante y y1 y2 applichiamo le formu-le di prostaferesi: y y1 y2 A cos (2πf1t) A cos (2πf2t) A[cos (2πf1t) cos (2πf2t)] 2A cos cos 2A cos2π t cos2π t L’ultima formula descrive l’andamento in funzione del tempo di un’onda in cui: • la frequenza (f1 f2)/2 è la media aritmetica delle frequenze delle due onde; • l’ampiezza dipende dal tempo secondo la legge: 2A cos2π t Poiché il volume del suono è nullo quando questo termine si annulla, e ciò avviene due volte ogni ciclo, la frequenza dei battimenti è: 2 f1 f2 In altri termini, la frequenza dei battimenti è uguale alla differenza tra le frequen-ze dei due suoni.
Nella situazione rappresentata nella figura 12.33 la frequenza dei battimenti è 2 Hz,(cioè 440 Hz 438 Hz),quindi un osservatore sente due battimenti al secondo.
La figura 12.34 mostra i grafici della pressione dell’onda sonora in un punto fissa-to in funzione del tempo per un’onda con frequenza di 10 Hz, per un’onda con fre-quenza di 12 Hz e per l’onda che risulta dalla loro sovrapposizione. Queste due fre-quenze sono state scelte solo per semplificare la rappresentazione grafica, perché in realtà sono minori del limite di udibilità e quindi un osservatore non sente alcun battimento. Comunque le onde con frequenze minori o maggiori dei limiti di udi-bilità si comportano esattamente come quelle udibili e quindi la descrizione del fenomeno dei battimenti è la stessa.
I due grafici disegnati in blu rappresentano le variazioni di pressione in un inter-vallo di tempo di 1 s per le due onde emesse dai due diapason. Il grafico disegnato in rosso rappresenta le variazioni di pressione della somma delle due onde calcola-ta applicando il principio di sovrapposizione. Osserviamo che l’ampiezza della pres-sione nel grafico in rosso non è costante, ma varia periodicamente da un valore mas-simo a un valore minimo. Quando queste variazioni di pressione raggiungono l’o-recchio di un osservatore e la loro frequenza appartiene all’intervallo di udibilità, esse producono un suono forte quando l’ampiezza della pressione è massima e un Tempo Tempo Tempo Debole Forte Debole Forte 1 1 2 3 4 5 6 7 8 9 10 11 12 2 3 4 5 6 7 8 9 10 f1 f2 2 f1 f2 2 f1 f2 2 f1 f2 2 2πf1t 2πf2t 2 2πf1t 2πf2t 2 CAPITOLO 12 Le onde e il suono 351 Figura 12.34 La sovrapposizione di due onde sonore con frequenze di 10 Hz e 12 Hz dà origine a un’onda sonora in cui sono presenti battimenti con una frequenza di 2 Hz. I grafici rappresentano come variano nel tempo le pressioni delle onde sonore che si sovrappongono (in blu) e la pressione dell’onda che risulta dalla sovrapposizione (in rosso).
L’intervallo di tempo tra due linee verticali tratteggiate è 1 secondo.
suono debole quando l’ampiezza della pressione è minima. Nel grafico in rosso si possono osservare due variazioni da debole a forte, cioè due battimenti in 1 s. Quin-di la frequenza dei battimenti è 2 Hz, pari alla differenza delle frequenze delle due onde sonore (12 Hz 10 Hz 2 Hz).
Il fenomeno dei battimenti è spesso usato dai musicisti per accordare i loro stru-menti. Per esempio, per accordare una corda del suo strumento, un chitarrista la piz-zica e contemporaneamente fa suonare un diapason o un altro strumento di cui sa che la frequenza è quella corretta e continua a regolare la tensione della corda fino a quando non sente più battimenti. Infatti, se non ci sono battimenti, le frequenze emesse dalla corda e dal diapason sono uguali.
12.10 Onde stazionarie Un importante fenomeno provocato dall’interferenza è quello delle onde staziona-rie, che possono essere prodotte sia dalle onde trasversali sia dalle onde sonore lon-gitudinali.
Modi normali La figura 12.35 mostra alcune delle caratteristiche principali delle onde stazionarie trasversali. In questa figura l’estremo sinistro di ciascuna corda viene fatto vibrare avanti e indietro,mentre l’estremo destro è fissato a una parete.Le regioni della corda si muovono così velocemente che le fotografie delle forme della corda sembrano sfuo-cate.Queste forme sono chiamate modi normali delle onde stazionarie trasversali.
Possiamo osservare che in queste forme sono presenti alcuni punti particolari, chiamati nodi e ventri. I nodi sono i punti in cui non c’è alcuna vibrazione, mentre i ventri sono i punti in cui l’ampiezza della vibrazione è massima. A destra di cia-scuna fotografia c’è un grafico che è sovrapponibile alla forma dell’onda e aiuta a visualizzare il moto della corda mentre vibra in uno dei suoi modi normali. I dise-gni rappresentano «fotografie» delle forme della corda in istanti diversi e mettono in evidenza l’ampiezza massima della vibrazione che si verifica in un ventre con pal-lini rossi attaccati alla corda.
CAPITOLO 12 Le onde e il suono 352 Fisica quotidiana Accordare uno strumento musicale 1 a armonica (frequenza fondamentale) 2a armonica 3a armonica f1 Frequenza 3f1 2f1 Ventri Nodi A B C Figura 12.35 Quando una corda fissata ai due estremi vibra a ben precise frequenze, essa diventa sede di onde stazionarie trasversali, come le tre onde mostrate nelle fotografie a sinistra. I disegni a fianco di ciascuna fotografia mostrano le forme assunte dalla corda nei tre casi e i pallini rossi attaccati alla corda mettono in evidenza le ampiezze massime delle vibrazioni che si verificano nei ventri delle onde.
© Richard Megna/Fundamental Photographs Ciascuno dei modi normali di un’onda stazionaria si verifica per un valore ben pre-ciso della frequenza della vibrazione. Queste frequenze formano una serie, chiamata serie armonica. Come mostra la figura 12.35, la frequenza minore della serie, chia-mata prima armonica o frequenza fondamentale e indicata con f1, corrisponde al modo normale in cui la corda ha due nodi e un ventre, mentre le frequenze suc-cessive sono multipli interi della frequenza fondamentale (2f1, 3f1, 4f1, ...) e sono chiamate rispettivamente seconda, terza, quarta, ... armonica. Il numero dell’armo-nica (prima, seconda, terza ecc.) corrisponde al numero di ventri nel modo norma-le dell’onda stazionaria.
L’origine delle onde stazionarie su una corda Le onde stazionarie si formano perché sulla corda si propagano onde identiche che viaggiano in versi opposti e si sommano come previsto dal principio di sovrapposi-zione. Il motivo per cui queste onde sono chiamate stazionarie è che non si propa-gano, né in un verso né nell’altro, come invece fanno le onde che le producono, e al passare del tempo mantengono le loro caratteristiche invariate.
La figura 12.36 spiega il motivo per cui su una corda in cui si formano onde sta-zionarie ci sono onde che viaggiano in versi opposti. Nel primo grafico dall’alto è rappresentato per semplicità solo mezzo ciclo di un’onda che viaggia verso la pare-te a destra. Quando questo mezzo ciclo raggiunge la parete, la corda esercita una spinta verso l’alto sulla parete. Di conseguenza, per la terza legge di Newton, la pare-te esercita una spinta verso il basso sulla corda e in questo modo produce un mezzo ciclo che viaggia sulla corda verso sinistra. In altre parole, l’onda viene riflessa dalla parete. Quando l’onda riflessa arriva al punto di origine, viene riflessa di nuovo, que-sta volta per effetto della vibrazione impartita alla corda dalla mano. Per vibrazioni di ampiezza piccola la mano rimane praticamente ferma e riflette l’onda come se fos-se una parete fissa. Queste continue riflessioni ai due estremi della corda danno ori-gine a un grande numero di cicli che viaggiano in versi opposti sulla corda.
Quando si forma un nuovo ciclo nell’estremo della corda mantenuto in vibrazio-ne dalla mano, arrivano alla mano i cicli precedenti che erano stati riflessi dalla parete. Tuttavia, se non si riesce a far vibrare la corda con la frequenza corretta, i cicli nuovi e quelli precedenti tendono a cancellarsi a vicenda quando si sovrap-pongono e non si formano le onde stazionarie. Pensiamo, per esempio, a che cosa succede se stiamo spingendo una persona su un’altalena e diamo ogni spinta nel momento giusto per rinforzare ogni volta l’ampiezza dell’oscillazione dell’altale-na. Questo è quello che succede quando le vibrazioni della mano hanno la fre-quenza giusta e ogni nuovo ciclo provoca la formazione di un’onda stazionaria di grande ampiezza.
Calcoliamo quale deve essere la frequenza delle vibrazioni impartite della mano per ottenere questo risultato. Supponiamo che la lunghezza della corda sia L e che il suo estremo sinistro venga fatto vibrare a una frequenza f1. L’intervallo di tempo necessario per produrre un nuovo ciclo è il periodo T dell’onda, dove T 1/f1. Inve-ce il tempo impiegato da un ciclo per viaggiare dalla mano alla parete e dalla pare-te alla mano, cioè per percorrere una distanza 2L, è 2L/v, dove v è la velocità del-CAPITOLO 12 Le onde e il suono 353 Figura 12.36 Quando il mezzo ciclo d’onda che viaggia verso la parete a destra viene riflesso dalla parete, diventa un mezzo ciclo d’onda che viaggia verso sinistra ed è capovolto.
FILM Le armoniche l’onda. I cicli nuovi rinforzano quelli precedenti quando questi due intervalli di tem-po sono uguali, cioè quando 1/f1 2L/v. Perciò, per generare un’onda stazionaria occorre far vibrare la corda con una frequenza: f1 I continui rinforzi tra i nuovi cicli prodotti dalla mano e i cicli riflessi portano alla formazione di un’onda stazionaria di grande ampiezza sulla corda, anche quando l’ampiezza delle vibrazioni impartite dalla mano è piccola.
Frequenze dei modi normali La frequenza minima a cui si deve far vibrare una corda per generare un’onda sta-zionaria è chiamata frequenza naturale della corda. Ogni corda ha però una serie di modi normali, a ciascuno dei quali corrisponde una ben precisa frequenza.
Questa serie è dovuta al fatto che, per generare un’onda stazionaria, non è neces-sario che ogni nuovo ciclo venga rinforzato da un ciclo che ritorna al punto di ori-gine. Per esempio, si può avere un rinforzo anche un nuovo ciclo sì e uno no, come avviene quando la corda viene fatta vibrare con la frequenza 2f1, oppure ogni tre nuovi cicli, come avviene quando la corda viene fatta vibrare con la frequenza 3f1.
Lo stesso ragionamento vale per qualunque frequenza fn nf1, dove n è un nume-ro naturale.
Di conseguenza, la serie di frequenze che danno luogo alla formazione di onde stazionarie in una corda fissata a entrambi i suoi estremi è Corda fissata a entrambi gli estremi fn n n 1, 2, 3, 4,… (12.12) Le onde stazionarie svolgono un ruolo molto importante nel modo in cui molti stru-menti musicali producono i loro suoni. Per esempio, una corda di chitarra è tesa tra due estremi fissi e quando viene pizzicata le sue vibrazioni sono quelle della serie di frequenze date dall’equazione (12.12).
v 2L v 2L CAPITOLO 12 Le onde e il suono 354 ESEMPIO 7 Onde stazionarie Suonare una chitarra Quando è pizzicata, la corda più pesante di una chitarra produce la nota mi. Un chitarrista vuole che la corda emetta il mi dell’ottava superiore. Per ottenere que-sto risultato deve premere il tasto giusto prima di pizzicare la corda (figura 12.37B).
Calcola la distanza L tra il tasto su cui premere la corda e il ponte della chi-tarra.
Ragionamento e soluzione La frequenza fondamentale è: f1 Per ottenere l’ottava superiore, cioè una frequenza doppia, con una corda di lun-ghezza x, deve essere: 2 da cui segue: x Quindi la lunghezza della corda oscillante deve essere L/2, cioè la metà della lun-ghezza totale della corda.
L 2 v 2x v 2L v 2L Figura 12.37 La figura mostra le onde stazionarie (in blu) che si formano su una corda di chitarra in diverse condizioni.
L L A B CAPITOLO 12 Le onde e il suono 355 IL MODELLO (decibel prodotti dall’altoparlante nel punto in cui è lo spettatore) 10 [logaritmo in base 10 di (intensità del-l’onda sonora nel punto in cui è lo spettatore / intensità minima udibile dall’orecchio umano)] I NUMERI Intensità dell’onda sonora prodotta dall’altoparlante nel punto in cui è lo spettatore (potenza sonora totale emessa dall’altoparlante) / (area attraverso la quale si distribuisce il segnale sonoro in cor-rispondenza del punto occupato dallo spettatore) (1 102 W) / (area della semisfera centrata sull’altopar-lante e passante per il punto in cui è lo spettatore) (1 102 W) / [2 (4m)2] (1 102 W) / (1 102 m2) 1 W/m2 Intensità minima udibile dall’orecchio umano 1,0 1012 W/m2 IL RISULTATO Intensità dell’onda sonora 10 log10 ( ) 10 log10 1012 L’ordine di grandezza è: Ascoltare la musica a 120 dB può provocare dolore e dan-ni permanenti all’udito.
Un paragone Se gli altoparlanti vicini fossero due invece di uno, come spesso accade, lo spettatore sarebbe esposto all’intensità sonora di soli 3 dB in più. Infatti 10 log10 ( ) dB 10 log10 (2 1012) dB 10 (log10 2 log10 1012) dB 10 (0,30 12) dB (3 120) dB 2 1 W/m2 1012 W/m2 102 dB 120 dB 1 W/m2 1012 W/m2 A quanti decibel è esposto uno spettatore della prima fila?
Per calcolare a quanti decibel è sottoposto uno spettatore che si trova a 4 metri da un altoparlante, bisogna moltiplicare per 10 il logaritmo in base 10 del rapporto fra l’intensità dell’onda nel punto in cui è lo spettatore e l’intensità minima udibile dall’orecchio umano.
L’ordine di grandezza Qual è la forza che si esercita sui tuoi timpani quando ascolti l’iPod a tutto volume?
IL MODELLO (forza esercitata sul timpano) (pressione dell’onda sonora) (superficie del timpano) I NUMERI Pressione dell’onda sonora (minima pressione acustica udibile) 10 (2 105 Pa) 10 11 Pa Superficie del timpano 9 105 m2 IL RISULTATO Forza esercitata dall’iPod .................. N Le fonti Pressione minima rilevabile dall’orecchio umano: Stanford University Center for Computer Research in Music and Acoustics (www-ccrma.stanford.edu/~jos/mdft/DB_SPL.html) Decibel prodotti dall’iPod a volume massimo: UAB Health System (www.health.uab.edu/17730/) Superficie del timpano: Anatomia del Gray, Zanichelli 115 20 (decibel prodotti dall’iPod a volume massimo) 20 Stima l’ordine di grandezza I concetti fondamentali 356 1. La natura delle onde Un’onda è una perturbazione che si propaga nello spazio trasportando energia ma non materia. In un’onda trasversale la direzione della perturbazione è perpendico-lare alla direzione di propagazione dell’onda. In un’onda longitudinale la direzione della perturbazione è parallela alla direzione di propagazione dell’onda.
2. Onde periodiche Onda periodica Un’onda periodica è una perturbazione formata da cicli che si ripetono uguali a se stessi per tutto il tempo in cui la sorgente della perturbazione continua a produrla.
L’ampiezza dell’onda è lo spostamento massimo di una particella del mezzo in cui si propaga l’onda dalla sua posizione di equilibrio. La lunghezza d’onda λ è la distan-za che l’onda percorre mentre compie un ciclo completo ed è uguale alla distanza tra due punti equivalenti successivi, come due creste.
Periodo e frequenza Il periodo T è l’intervallo di tempo impiegato per compiere un intero ciclo ed è ugua-le al tempo impiegato dall’onda per percorrere una distanza uguale a una lunghez-za d’onda. La frequenza f (in hertz) è il numero di cicli che avvengono in 1 s ed è uguale all’inverso del periodo (in secondi): f Velocità di propagazione La velocità di propagazione v di un’onda è legata alla sua lunghezza d’onda e alla sua frequenza dalla relazione: v f λ (12.1) La velocità di propagazione di un’onda dipende dalle caratteristiche del mezzo in cui si propaga.Per un’onda trasversale che si propaga in una corda la velocità è legata alla tensione T della corda e alla sua massa per unità di lunghezza m/L dalla relazione: v (12.2) 3. La descrizione matematica di un’onda L’equazione che descrive un’onda di ampiezza A, frequenza f e lunghezza d’onda λ che si propaga lungo la direzione x è: Onda che si propaga verso x y A sen 2πft (12.3) Onda che si propaga verso x y A sen 2πft (12.4) L’angolo 2πft è detto fase dell’onda e si deve misurare in radianti.
4. La natura del suono Onde sonore Il suono è formato da onde longitudinali, generate da un oggetto che vibra, che si propagano in un mezzo materiale e non si propagano nel vuoto. Ogni ciclo di un’on-da sonora è formato da una compressione e una rarefazione.
Suoni puri e suoni complessi Un suono si dice puro quando le particelle investite dall’onda sonora oscillano con moto armonico: la frequenza di oscillazione delle particelle è la frequenza del suono puro. Quando oscillano di moto periodico ma non armonico, il suono si dice com-plesso: anche in questo caso, si può individuare una frequenza che caratterizza il suo-no, detta frequenza fondamentale.
2πx λ 2πx λ 2πx λ Equazione di un’onda periodica che si propaga nella direzione x T m/L Velocità di un’onda trasversale su una corda 1 T Ampiezza e lunghezza d’onda Onde trasversali e onde longitudinali CAPITOLO 12 Le onde e il suono 357 Infrasuoni e ultrasuoni Le onde sonore con frequenza minore di 20 Hz sono chiamate infrasuoni, mentre quelle con frequenza maggiore di 20 kHz sono chiamate ultrasuoni.
Suoni alti e suoni bassi Il nostro cervello interpreta le diverse frequenze delle onde che arrivano all’orecchio come suoni di altezza diversa: i suoni alti (o acuti) sono quelli con frequenza mag-giore (per esempio, sono acuti i suoni emessi da un ottavino), mentre i suoni bassi (o gravi) sono quelli con frequenza minore (per esempio, sono gravi i suoni emessi da una tuba).
Ampiezza della pressione sonora L’ampiezza della pressione di un’onda sonora è il valore massimo della differenza di pressione tra una regione di compressione e la pressione normale del mezzo in cui l’onda si propaga. Dall’ampiezza della pressione dipende una caratteristica soggetti-va del suono che è l’intensità sonora: quanto maggiore è l’ampiezza della pressione, tanto più forte è il suono percepito.
5. L’intensità del suono Intensità di un suono L’intensità di un suono I è il rapporto tra la potenza sonora media – P che attraversa perpendicolarmente una superficie e l’area A della superficie: I (12.5) L’unità di misura SI dell’intensità del suono è il watt al metro quadrato (W/m2).
Soglia minima udibile L’intensità minima percepibile da un orecchio umano è chiamata soglia minima udi-bile e per un suono con una frequenza di 1 kHz vale circa 1 10 12 W/m2.
Onda sferica uniforme Quando una sorgente emette un’onda sonora che si propaga uniformemente nello spazio e non sono presenti ostacoli che possano riflettere l’onda, l’intensità del suo-no in un punto è inversamente proporzionale al quadrato della distanza del punto dalla sorgente: I (12.6) Livello di intensità sonora L’intensità con cui percepiamo un suono di intensità I si misura mediante il livello di intensità sonora β, così definito: β 10 log10 dove I0 1 1012 W/m2 è la minima intensità sonora udibile. Nel Sistema Interna-zionale il livello di intensità sonora è espresso in decibel (dB).
6. L’effetto Doppler Effetto Doppler L’effetto Doppler è la variazione di frequenza tra il suono emesso da una sorgente e quello percepito da un osservatore, perché la sorgente o il ricevitore o entrambi si muovono con velocità diverse rispetto al mezzo in cui si propaga l’onda sonora. Se la sorgente e il ricevitore si muovono con velocità rispettivamente vs e vr e il mezzo è stazionario, la frequenza fr del suono percepito dal ricevitore è: fr fs (12.11) dove fs è la frequenza del suono emesso dalla sorgente e v è la velocità del suono.
Nel numeratore si deve usare il segno più quando il ricevitore si muove verso la sor-gente e il segno meno quando il ricevitore si allontana dalla sorgente. Nel denomi-natore si deve usare il segno meno quando la sorgente si muove verso il ricevitore e il segno più quando la sorgente si allontana dal ricevitore.
1 — vr v 1 — vs v I I0 P 4πr2 – P A CAPITOLO 12 Le onde e il suono 358 7. Il principio di sovrapposizione Principio di sovrapposizione Il principio di sovrapposizione afferma che, quando due o più onde sono presenti contemporaneamente in uno stesso punto, la perturbazione in quel punto è la som-ma delle perturbazioni prodotte dalle singole onde.
8. Interferenza e diffrazione di onde sonore Quando due onde si combinano in modo che una compressione si sovrapponga sem-pre esattamente a una compressione e una rarefazione si sovrapponga sempre esat-tamente a una rarefazione, si dice che sono in fase (o in concordanza di fase) e che danno luogo a interferenza costruttiva.
Quando due onde si combinano in modo che una compressione si sovrapponga sem-pre esattamente a una rarefazione, si dice che sono in opposizione di fase e che dan-no luogo a interferenza distruttiva.
Se la differenza tra i cammini percorsi dalle onde emesse da due sorgenti sonore che vibrano in fase per arrivare al punto in cui si sovrappongono è pari a zero o a un numero intero (1, 2, 3,…) di lunghezze d’onda si ha il fenomeno di interferen-za costruttiva. Una differenza tra i cammini percorsi pari a mezza lunghezza d’on-da o a un numero intero di lunghezze d’onda più mezza lunghezza d’onda , 1 , 2 , 3 , … produce un fenomeno di interferenza distruttiva.
Diffrazione Quando incontra un ostacolo oppure i bordi di una fenditura, un’onda devia dalla sua direzione di propagazione e prosegue al di là dell’ostacolo o della fenditura; que-sto fenomeno è detto diffrazione.
9. Battimenti Battimenti I battimenti sono variazioni periodiche dell’ampiezza provocate dall’interferenza di due onde con frequenze leggermente diverse. Se le onde sono onde sonore, esse ven-gono percepite come un suono unico la cui intensità varia con una frequenza ugua-le alla differenza tra le frequenze delle due onde.
10. Onde stazionarie Onde stazionarie Un’onda stazionaria è il tipo di perturbazione che si verifica quando si sovrappon-gono due onde che hanno la stessa ampiezza, la stessa frequenza e versi di propaga-zione opposti.
Nodi e ventri I punti di un’onda stazionaria in cui l’ampiezza della vibrazione è nulla si chiamano nodi, mentre quelli in cui l’ampiezza della vibrazione è massima si chiamano ventri.
Serie armonica Le onde stazionarie possono formarsi solo quando le frequenze della vibrazione han-no valori ben precisi che fanno parte di una serie chiamata serie armonica. Le fre-quenze di questa serie (f1, 2f1, 3f1 ecc.) sono chiamate armoniche. La frequenza più bassa f1 è chiamata prima armonica o frequenza fondamentale, la frequenza succes-siva 2f1 è chiamata seconda armonica e così via.
Per una corda di lunghezza L fissata a entrambi gli estremi le frequenze della serie armonica sono date da: fn n n 1, 2, 3, 4,… (12.12) dove v è la velocità del suono e n è un numero naturale.
v 2L Frequenze della serie armonica per una corda 1 2 1 2 1 2 1 2 Condizioni per l’interferenza costruttiva e distruttiva Interferenza costruttiva e interferenza distruttiva CAPITOLO 12 Le onde e il suono 359 In una molla elicoidale si propaga un’onda alla velocità di 5 m/s. Questo significa che una spira della molla si muove di 5 m in 1 s?
Un’onda sonora si propaga in un mezzo. Esistono parti-celle del mezzo che rimangono sempre in quiete?
Una sorgente emette in modo uniforme un suono che, senza alcuna riflessione, è in parte intercettato da una superficie piana. L’intensità sonora è la stessa in tutti i punti della superficie? 3 2 1 Quando un’auto è ferma, il suo clacson emette un suo-no con una frequenza di 600 Hz. Una persona in mezzo alla strada sente il clacson a una frequenza di 580 Hz.
Deve affrettarsi a raggiungere il marciapiede?
Il principio di sovrapposizione implica che due onde sonore che passano nello stesso istante per lo stesso pun-to danno sempre luogo a un suono più forte?
Se la tensione della corda di una chitarra viene raddop-piata, anche la frequenza raddoppia? In caso negativo, di quale fattore cambia? Aumenta o diminuisce?
6 5 4 Domande Quale fra le seguenti è un’onda longitudinale?
a Onda sonora nell’aria.
b Onda luminosa nell’aria.
c Onda radio nell’aria.
d Onda sull’acqua.
La velocità del suono in un dato metallo è 3,00 103 m/s.
Il grafico mostra l’ampiezza in metri in funzione del tempo in millisecondi di un’onda che viaggia nel metal-lo. Qual è la sua lunghezza d’onda?
a 0,5 m c 4,0 m b 1,5 m d 6,0 m Un’onda ha una frequenza di 58 Hz e si propaga alla velocità di 31 m/s. Qual è la sua lunghezza d’onda?
a 0,29 m b 0,53 m c 1,9 m d 3,5 m Nella figura seguente è rappresentata un’onda che si muove a 10,0 cm/s. Qual è l’ampiezza dell’onda?
a 2 cm b 4 cm c 6 cm d 12 cm y (cm) 6 4 2 0 2 4 6 1,5 4,5 7,5 10,5 13,5 16,5 19,5 22,5 x (cm) 4 3 y (m) 1 1,0 2,0 3,0 4,0 1 0 t (ms) 2 1 Qual è la frequenza dell’onda rappresentata nel grafico del quesito precedente?
a 1,7 Hz b 1,3 Hz c 1,1 Hz d 0,9 Hz Un’onda si propaga in un mezzo con la legge: y (2 m) sen [(π s1) t (π m1) x].
La sua frequenza f e la sua lunghezza d’onda λ sono rispettivamente: a f π Hz λ π m b f π Hz λ 1/π m c f 0,5 Hz λ 2 m d f 2 Hz λ 0,5 m Un campanello elettrico è posto sotto una campana di vetro. Mentre sta suonando, l’aria viene tolta lentamen-te dalla campana. Che cosa sente un ascoltatore?
a L’intensità del suono rimane costante.
b L’intensità del suono diminuisce progressivamente.
c La frequenza del suono aumenta progressivamente.
d La frequenza del suono diminuisce progressiva-mente.
Un suono si dice puro quando: a è emesso da uno strumento musicale.
b le particelle investite dall’onda sonora oscillano con moto armonico.
c le particelle investite dall’onda sonora oscillano con moto periodico.
d le particelle investite dall’onda sonora si muovono in modo uniforme.
L’intensità di un’onda sferica a 4,0 m dalla sorgente è 120 W/m2. Qual è l’intensità a 9 m dalla sorgente?
a 11 W/m2 b 24 W/m2 c 53 W/m2 d 80 W/m2 9 8 7 6 5 Test Esercizi CAPITOLO 12 Le onde e il suono 360 Un treno transita in una stazione a velocità costante. Sul treno un flautista emette un la a 440 Hz. Quando il treno si allontana,il suono emesso viene percepito dal caposta-zione come sol a 392 Hz.Qual è la velocità del treno?
a 7,3 m/s b 12 m/s c 26 m/s d 42 m/s Un automobilista si avvicina a 35 m/s a un segnale acu-stico di 220 Hz. Con quale frequenza lo percepisce?
(Velocità del suono: 343 m/s.) a 198 Hz b 220 Hz c 242 Hz d 282 Hz Due impulsi della stessa ampiezza viaggiano su una cor-da l’uno verso l’altro a 1,0 m/s. La figura mostra la loro posizione all’istante t 0 s. Qual è la forma della corda all’istante t 2,0 s?
a b c d y 2 m x 2 m y 2 m x 2 m y 2 m x 2 m y 2 m x x 2 m y 2 m 2 m 12 11 10 1. La natura delle onde 2. Onde periodiche La luce è un’onda elettromagnetica che viaggia a una velocità di 3,00 108 m/s. La frequenza luminosa a cui l’occhio umano è più sensibile è quella della luce giallo-verde, che ha una lunghezza d’onda di 5,45 107 m.
Qual è la frequenza di questa onda luminosa?
Una persona su una barca ferma nel mare osserva che, dopo il passaggio della cresta di un’onda, passano altre 5 creste in 50 secondi. La distanza fra due creste suc-cessive è 32 m.
Determina,se possibile,il periodo,la frequenza,la lun-ghezza d’onda, la velocità e l’ampiezza dell’onda.
Un’onda longitudinale che ha una frequenza di 3,0 Hz impiega 1,7 s per attraversare una molla lunga 2,5 m.
Determina la sua lunghezza d’onda.
Supponi che la mano che fa muovere su e giù l’estremo di una molla faccia compiere all’estremo due oscillazio-ni complete in 1 s e che la velocità di propagazione del-4 3 2 1 l’onda così generata sia di 0,50 m/s.
Calcola la distanza tra due creste adiacenti dell’onda.
The speed of a transverse wave on a string is 450 m/s, and the wavelength is 0.18 m.The amplitude of the wave is 2.0 mm.
How much time is required for a particle of the string to move through a total distance of 1.0 km?
Una corda ha una massa di 5,0 103 kg ed è sottoposta a una tensione di 180 N. Un’onda trasversale che si pro-paga lungo questa corda ha una frequenza di 260 Hz e una lunghezza d’onda di 0,60 m.
Qual è la lunghezza della corda?
Una persona che fa sci d’acqua viaggia a una velocità di 12,0 m/s. Quando viaggia nella stessa direzione e nello stesso verso di un’onda, viene sollevata dalle creste del-l’onda una volta ogni 0,600 s. Quando viaggia nella stes-sa direzione ma in verso opposto viene sollevata dalle creste dell’onda una volta ogni 0,500 s. La sua velocità è maggiore di quella dell’onda.
Calcola la velocità e la lunghezza d’onda dell’onda.
7 6 5 Un ascoltatore dista 19 m da una cassa acustica e 15 m dall’altra. Le due casse emettono in fase lo stesso suono puro di frequenza f. La velocità del suono è 344 m/s.
Qual è il più piccolo valore di f per il quale l’ascoltato-re percepisce l’intensità maggiore?
a L’intensità non dipende da f.
b Non si può stabilire perché non è nota l’ampiezza.
c f 43 Hz d f 86 Hz Il fenomeno dei battimenti è dovuto alla: a sovrapposizione di onde che viaggiano con velocità leggermente diverse.
b sovrapposizione di onde che viaggiano in versi oppo-sti.
c sovrapposizione di onde che hanno ampiezze legger-mente diverse.
d sovrapposizione di onde che hanno frequenze leg-germente diverse.
Su una corda si formano onde stazionarie per effetto della: a sovrapposizione di onde identiche che viaggiano con velocità diverse.
b sovrapposizione di onde identiche che viaggiano in versi opposti.
c sovrapposizione di onde che hanno stessa velocità ma ampiezze diverse.
d sovrapposizione di onde che hanno diversa velocità ma stessa ampiezza.
15 14 13 Problemi CAPITOLO 12 Le onde e il suono 361 esercizi Un’onda trasversale si propaga con una velocità di 300 m/s su una corda orizzontale.
Se la tensione della corda viene aumentata di quattro volte, quale diventa la velocità di propagazione del-l’onda?
Due fili metallici entrambi di lunghezza uguale a 50,0 m e con la stessa densità lineare di 0,020 kg/m sono paral-leli tra loro e disposti uno sopra l’altro. La tensione nel filo A è pari a 6,00 102 N, mentre quella nel filo B è pari a 3,00 102 N. Due impulsi trasversali sono genera-ti contemporaneamente all’estremo sinistro del filo A e all’estremo destro del filo B.
Quanto tempo trascorre prima che i due impulsi pas-sino uno sopra l’altro?
Per misurare l’accelerazione di gravità su un pianeta lontano un astronauta appende una sfera con una mas-sa di 0,55 kg a un capo di un filo metallico lungo 0,95 m la cui densità lineare è di 1,2 104 kg/m. L’astronauta misura il tempo impiegato da un impulso trasversale a propagarsi lungo il filo e ottiene un valore di 0,016 s.
La massa del filo è trascurabile rispetto a quella della palla.
Qual è l’accelerazione di gravità su quel pianeta?
3. La descrizione matematica di un’onda Lo spostamento in metri di un’onda è: y 0,26 m sen [(πs1) t (3,7 πm1) x].
dove t è in secondi e x in metri.
Qual è il verso di propagazione dell’onda?
Qual è lo spostamento quando t 38 s e x 13 m?
Un’onda si propaga nel verso positivo delle x e ha ampiezza di 0,35 m, velocità 5,2 m/s e frequenza 14 Hz.
Scrivi l’espressione matematica dell’onda.
I due grafici sono relativi a un’onda che si propaga nel-la direzione x.
Scrivi l’espressione matematica dell’onda.
Un’onda trasversale si propaga su una corda. Lo spo-stamento y delle particelle della corda dalla posizione di equilibrio è dato dalla legge: y (0,021 m) sen [(25 s1) t (2,0 m1) x].
L’angolo di fase [(25 s1) t (2,0 m1) x] è in radianti, t in secondi e x in metri. La densità lineare della corda è 1,6 102 kg/m.
Calcola la tensione della corda.
14 0,020 0,060 0,040 – 0,010 + 0,010 0,080 t = 0 s x (m) y (m) 0,10 0,30 0,20 – 0,010 + 0,010 0,40 x = 0 m t (s) y (m) 13 12 11 10 9 8 4. La natura del suono 5. L’intensità del suono La distanza fra un altoparlante e l’orecchio sinistro di un ascoltatore è 2,70 m. Nell’aria della stanza il suono si propaga a 343 m/s e ha una frequenza di 523 Hz.
Calcola il tempo che il suono impiega a raggiungere l’ascoltatore.
Calcola il numero di lunghezze d’onda presenti fra l’altoparlante e l’ascoltatore.
Come mostra la figura, si può ottenere una sirena sof-fiando un getto d’aria attraverso 20 fori praticati alla stessa distanza tra loro sul bordo di un disco rotante. Il tempo impiegato da ciascun foro per passare davanti al getto d’aria è il periodo del suono emesso dalla sirena.
Se si vuole che questo suono abbia una frequenza di 2200 Hz, quale deve essere la velocità angolare ω (in rad/s) del disco rotante?
Un terremoto genera due tipi di onde sismiche che si propagano attraverso il globo terrestre: le onde prima-rie, o onde P, sono onde longitudinali che si propagano più velocemente delle onde secondarie, o onde S, che sono invece onde trasversali. Le onde P hanno una velo-cità di circa 8,0 km/s, mentre le onde S hanno una velo-cità di circa 4,5 km/s. Un sismografo che si trova a una certa distanza dall’ipocentro di un terremoto (cioè dal punto di origine delle onde) registra l’arrivo delle onde P e dopo 78 s l’arrivo delle onde S.
Supponendo che le onde sismiche si propaghino in linea retta, qual è la distanza tra il sismografo e l’ipo-centro del terremoto?
La superficie media dell’orecchio di un adulto è di 2,1 103 m2. L’intensità del suono che giunge all’orec-chio durante una conversazione normale è di circa 3,2 106 W/m2.
Calcola la potenza sonora che giunge sull’orecchio quando il suono arriva in direzione perpendicolare a esso.
Alla distanza di 3,8 m, l’intensità del suono emesso da una sirena è 3,6 102 W/m2. Il suono si propaga unifor-memente in tutte le direzioni.
Calcola la potenza sonora totale emessa dalla sirena.
Supponi che un suono venga emesso uniformemente in tutte le direzioni da un altoparlante. L’intensità del suo-no in un punto distante 22 m dall’altoparlante è pari a 3,0 104 W/m2.
Qual è l’intensità del suono in un punto a una distan-za di 78 m dall’altoparlante?
20 19 18 17 Getto d'aria 16 15 La terapia con ultrasuoni viene spesso impiegata per accelerare la guarigione di un tendine lesionato. Il prin-cipio di funzionamento si basa sull’energia che gli ultra-suoni cedono ai tessuti che attraversano e che viene con-vertita in calore. Gli ultrasuoni escono da un foro cir-colare con il raggio di 1,8 cm e sono emessi da un ge-neratore che produce un’intensità sonora pari a 5,9 103 W/m2.
Quanto tempo impiega il generatore per emettere 4800 J di energia sonora?
In uno studio di registrazione il livello sonoro è 44 dB più basso dell’esterno. Un mattino l’intensità sonora nel-lo studio è 1,20 1010 W/m2.
Calcola l’intensità sonora esterna.
Un apparecchio acustico aumenta l’intensità sonora di 30 dB.
Calcola di quale fattore aumenta l’intensità sonora.
Quando una persona è esposta per 9,0 ore a un livello di intensità sonora di 90,0 dB riporta seri danni acusti-ci. L’area di un timpano è 9,0 105 m2.
Calcola quanta energia sonora incide sul timpano durante quell’intervallo di tempo.
6. L’effetto Doppler L’antifurto di un’automobile parcheggiata emette un suono di frequenza pari a 960 Hz. La velocità del suono è di 343 m/s. Avvicinandoti, rilevi che la frequenza è cambiata di 95 Hz.
Qual è la tua velocità?
Mentre stai andando in bicicletta ti allontani da una sor-gente sonora e la frequenza del suono che senti è mino-re dell’1% rispetto alla frequenza del suono emesso dal-la sorgente.
A quale velocità stai viaggiando?
Sei fermo a un semaforo e un’ambulanza si avvicina a 18 m/s. La sirena dell’ambulanza emette un suono con frequenza pari a 995 Hz. La velocità del suono nell’aria è 343 m/s.
Qual è la lunghezza d’onda del suono che senti?
Il pilota di un aereo ultraleggero sta volando a 39 m/s.
Un’aquila vola verso l’aereo a 18 m/s. Le velocità sono riferite al terreno. L’aquila emette un fischio a 340 Hz.
Calcola la frequenza rilevata dal pilota.
Una portaerei viaggia a una velocità di 13,0 m/s rispet-to all’acqua.Un aereo viene catapultato dal ponte di lan-cio della portaerei e ha una velocità di 67,0 m/s rispetto all’acqua. I suoi motori producono un suono di fre-quenza pari a 1550 Hz e la velocità del suono nell’aria è di 343 m/s.
Qual è la frequenza del suono percepito dall’equipag-gio della portaerei?
Due sottomarini si muovono uno verso l’altro viaggian-do sott’acqua. La velocità del sottomarino A è 12 m/s, 30 29 28 27 26 25 24 23 22 21 mentre quella del sottomarino B è 8 m/s. Il sottomarino A emette un’onda sonora con una frequenza di 1550 Hz e una velocità di 1552 m/s.
Qual è la frequenza dell’onda rilevata dal sottomari-no B?
Una parte dell’onda emessa dal sottomarino A viene riflessa dal sottomarino B e ritorna al sottomarino A.
Qual è la frequenza dell’onda riflessa rilevata dal sot-tomarino A?
7. Il principio di sovrapposizione 8. Interferenza e diffrazione di onde sonore Il grafico rappresenta una corda su cui si propagano due impulsi, diretti l’uno verso l’altro, che viaggiano con una velocità di modulo costante uguale a 1 cm/s nell’istante t 0 s.
Usando il principio di sovrapposizione disegna la for-ma degli impulsi negli istanti t 1 s, t 2 s, t 3 s e t 4 s.
Due altoparlanti posti uno dietro l’altro emettono due onde sonore identiche con frequenza 245 Hz.
Qual è la distanza minima a cui devono trovarsi i due altoparlanti perché un ascoltatore posto davanti a essi percepisca un’interferenza distruttiva? La velocità del suono è 343 m/s.
Due impulsi viaggiano l’uno verso l’altro con la stessa velocità di modulo 1 cm/s. La figura rappresenta le loro posizioni nell’istante t 0 s. Quando t 1 s, qual è l’am-piezza dell’impulso risultante in corrispondenza di: x 3 cm x 4 cm?
Supponi che la distanza tra i due altoparlanti A e B del-la figura 12.31 sia di 5,00 m e che gli altoparlanti siano in fase. Essi emettono due suoni identici con una fre-quenza di 125 Hz e la velocità del suono è 343 m/s.
Qual è la massima distanza possibile tra l’altoparlan-te B e l’ascoltatore C perché l’ascoltatore C rilevi un’interferenza distruttiva?
Supponi che la distanza tra i due altoparlanti A e B del-la figura 12.31 sia di 2,50 m. Gli altoparlanti oscillano a 429 Hz ma in opposizione di fase. La velocità del suono è 343 m/s.
35 34 1 cm 1 cm/s 1 cm/s 0 1 0 2 3 4 5 Distanza, cm 6 33 32 1 cm/s 1 cm/s 0 2 4 6 8 10 12 Distanza, cm 31 CAPITOLO 12 Le onde e il suono 362 L’osservatore C rileva un’interferenza costruttiva o distruttiva quando la sua distanza dall’altoparlante B è 1,15 m e 2,00 m?
La figura mostra due altoparlanti e un ascoltatore C. Gli altoparlanti emettono in fase un suono di 68,6 Hz. La velocità del suono è 343 m/s.
Calcola la minima distanza da A alla quale deve esse-re posto B perché l’ascoltatore non oda alcun suono.
Gli altoparlanti A e B emettono in concordanza di fase un segnale a 73,0 Hz. Essi sono orientati uno verso l’al-tro e distano 7,80 m. La velocità del suono è 343 m/s. Nel segmento che li congiunge esistono tre punti in cui si ha interferenza costruttiva.
Calcola le distanze di questi punti da A.
9. Battimenti 10. Onde stazionarie Due suoni puri sono emessi contemporaneamente. I gra-fici mostrano la variazione della pressione di ciascuno di essi in funzione del tempo.
Qual è la frequenza dei battimenti?
Due pianoforti emettono la stessa nota ma non sono accordati. Uno emette un suono di lunghezza d’onda 0,769 m e l’altro di 0,776 m. La velocità del suono è pari a 343 m/s.
Calcola il periodo dei battimenti.
La corda di una chitarra deve essere accordata. Suo-nando contemporaneamente la corda e un diapason a 440 Hz si odono battimenti di frequenza 3 Hz. Aumen-tando la tensione della corda, la sua frequenza aumen-ta e la frequenza dei battimenti diminuisce.
Qual era la frequenza originaria della corda?
In un violoncello la corda che ha la densità lineare mag-giore (1,56 102 kg/m) è quella del do. Questa corda produce una frequenza fondamentale di 65,4 Hz e la lun-ghezza del tratto di corda compreso tra i due estremi fis-si è 0,800 m.
Calcola la tensione della corda.
Una corda lunga 0,28 m è fissata a entrambi gli estremi.
La corda viene pizzicata e su di essa si forma un’onda 42 41 40 39 0,020 s Pressione Tempo Tempo Pressione 0,024 s 0 0 38 37 B 60,0° A 1,00 m C 36 stazionaria che vibra alla frequenza della seconda armo-nica. Le onde si muovono sulla corda a una velocità di 140 m/s.
Qual è la frequenza della seconda armonica?
La frequenza fondamentale di una corda fissata agli estremi è 256 Hz.
Quanto tempo impiega un’onda a percorrere la corda?
Una corda ha densità lineare di 8,5 103 kg/m ed è sot-toposta a una tensione di 280 N. La corda è lunga 1,8 m, è fissata agli estremi e oscilla come mostrato in figura.
Determina la velocità, la lunghezza d’onda e la fre-quenza delle onde che formano l’onda stazionaria.
Come mostra la figura, la lunghezza di una corda di una chitarra è 0,628 m. Per semplicità i tasti sono numera-ti. Un chitarrista può suonare tutte le note della scala su una sola corda perché le distanze tra un tasto e l’al-tro sono state progettate secondo la seguente regola: quando la corda viene premuta contro un tasto qua-lunque di numero i, la sua frequenza fondamentale di vibrazione è maggiore di un fattore pari alla radice dodicesima di 2 (12 2) rispetto alla frequenza fonda-mentale di vibrazione che ha quando viene premuta contro il tasto numero i 1.
Supponendo che la tensione della corda rimanga ugua-le per tutte le note, calcola la distanza tra la barretta 1 e la barretta 0 e tra la barretta 7 e la barretta 6.
PROBLEMI FINALI La densità lineare della corda del la di un violino è 7,8 104 kg/m. Un’onda sulla corda ha una frequenza di 440 Hz e una lunghezza d’onda di 65 cm.
Calcola la tensione della corda.
I 15 vagoni di un treno merci passano a velocità costan-te davanti ad un incrocio in 12,0 s. Ciascun vagone è lun-go 14,0 m.
Qual è la frequenza con cui passano i vagoni? Qual è la velocità del treno?
47 46 0 2 3 4 5 6 7 1 0,628 m 45 44 43 CAPITOLO 12 Le onde e il suono 363 esercizi Riprendi in considerazione i grafici del problema 13.
Calcola la velocità dell’onda.
Un aereo vola in direzione orizzontale, come mostra la figura. Quando arriva nel punto B che si trova sulla ver-ticale di un osservatore a terra, l’osservatore sente il rumore prodotto dall’aereo quando si trovava nel pun-to A. La velocità del suono è 343 m/s.
Se la velocità dell’aereo nel punto A è 164 m/s, qual è la sua velocità nel punto B, supponendo che l’aereo viaggi con accelerazione costante?
The average sound intensity inside a busy neigh-borhood restaurant is 3.2 105 W/m2.
How much energy goes into each ear (area 2.1 103 m2) during a one-hour meal?
Mentre vola, un cardellino emette un richiamo di fre-quenza 1250 Hz. Una naturalista registra il richiamo, che risulta essere di 1290 Hz.
Qual è la velocità con cui volava il cardellino? Espri-mila in percentuale della velocità del suono.
II suono prodotto dall’altoparlante rappresentato in figura ha una frequenza di 12 000 Hz e arriva al mi-crofono seguendo due cammini diversi: viaggia lungo il tubo di sinistra LXM, che ha una lunghezza fissa, e contemporaneamente viaggia lungo il tubo di destra LYM, la cui lunghezza può essere cambiata muovendo la parte di tubo scorrevole. Le onde sonore che seguo-no i due cammini interferiscono nel punto M. Quando la lunghezza del cammino LYM cambia, cambia anche l’intensità del suono che arriva nel punto M. Quando il tubo di destra viene tirato verso l’esterno di 0,020 m, l’intensità passa da un valore massimo a un valore minimo.
Trova la velocità a cui si propaga il suono nel gas con-tenuto nel tubo.
Due altoparlanti che vibrano in fase sono disposti nello stesso modo indicato nella figura 12.31 e anche l’ascol-tatore C si trova alla stessa distanza da essi. La velocità 53 Altoparlante L M Y X Microfono Parte di tubo scorrevole 52 51 50 A B 36,0° 49 48 del suono è 343 m/s. Gli altoparlanti emettono due suo-ni identici con la stessa frequenza.
Qual è la frequenza minima perché l’ascoltatore C oda un’interferenza distruttiva?
Una corda del la di un contrabbasso è accordata per vibrare a una frequenza fondamentale di 55,0 Hz.
Se la tensione della corda fosse aumentata di quattro volte, quale sarebbe la sua frequenza fondamentale di vibrazione?
Due ultrasuoni si combinano e danno luogo a battimen-ti che sono udibili dall’uomo (frequenza compresa fra 20 Hz e 20 000 Hz). La frequenza di uno di essi è 70 kHz.
Calcola la frequenza minore e la frequenza maggiore che può avere l’altro ultrasuono.
Le due corde rappresentate in figura hanno lunghezze e densità lineari diverse.Sono unite tra loro e tese in modo che la tensione in ciascuna di esse sia di 190,0 N. I due capi delle corde che non sono uniti tra loro sono fissati a due pareti.
Calcola la frequenza minima di vibrazione per la qua-le nelle due corde si formano onde stazionarie con un nodo nel punto di congiunzione.
QUESITI Illustra il fenomeno delle onde stazionarie e definisci i modi normali di oscillazione.
Descrivi e analizza il fenomeno dei battimenti.
Descrivi e analizza l’effetto Doppler quando la sorgen-te è ferma e il ricevitore in moto e viceversa. Confronta le frequenze registrate dal ricevitore.
Descrivi il fenomeno dell’interferenza e spiega quando si realizzano l’interferenza costruttiva e quella distruttiva.
Spiega che cos’è un’onda sonora e come è possibile generarla.
OLIMPIADI DELLA FISICA L’equazione di un’onda armonica è: y Y0 sen [k(x vt)] Se Y0 3 m, k 3π m1 e v 8 ms1, qual è la frequen-za dell’onda?
a 3,0 Hz b 7,2 Hz c 8,0 Hz d 12 Hz e 24 Hz (Gara di 1° livello edizione 2007) 1 5 4 3 2 1 6,00 102 kg/m 1,50 102 kg/m Nodo 3,75 m 1,25 m 56 55 54 CAPITOLO 12 Le onde e il suono 364 Una sorgente di onde acustiche di frequenza costante e un osservatore si muovono l’uno rispetto all’altra. La frequenza dell’onda misurata dall’osservatore cresce uniformemente col tempo. Questo accade perché rispet-to a un certo sistema di riferimento: a l’osservatore si muove a velocità costante verso la sorgente che è ferma.
b la sorgente si allontana a velocità costante dall’os-servatore che è fermo.
c l’osservatore si muove di moto uniformemente acce-lerato verso la sorgente ferma.
d la sorgente si allontana di moto uniformemente acce-lerato dall’osservatore fermo.
e l’osservatore percorre, a velocità costante, una cir-conferenza nel cui centro sta la sorgente.
(Gara di 1° livello edizione 2007) Il disegno schematizza un’onda d’acqua che si propaga alla velocità di 1 ms1 e mette in moto un tappo che com-pie 8 oscillazioni in 4 s.
Qual è la lunghezza d’onda?
a 0,25 m b 0,5 m c 1 m d 2 m e 4 m (Gara di 1° livello edizione 2006) In un determinato mezzo si propagano, nella stessa dire-zione, due impulsi che a un certo punto si sovrappongo-no. In quell’istante si ottiene l’impulso mostrato in figu-ra, dovuto alla loro sovrapposizione.
Fra i seguenti, quali possono essere gli impulsi iniziali?
(Gara di 1° livello edizione 2005) a b c d e 4 v Tappo 3 2 Un generatore di onde, collocato a 4 metri di distanza da un muro riflettente, produce onde stazionarie in un nastro, come mostrato nel disegno sotto.
Se la velocità dell’onda vale 10 ms1, qual è la sua fre-quenza?
a 0,4 Hz b 4 Hz c 5 Hz d 10 Hz e 40 Hz (Gara di 1° livello edizione 2004) Il grafico in figura riporta lo spostamento dalla loro posizione di equilibrio di particelle che, a un certo istan-te, risentono dell’effetto di un’onda che viaggia lungo l’asse x. P è un punto lungo il percorso fatto dall’onda.
Quale dei grafici seguenti rappresenta meglio lo sposta-mento in funzione del tempo di una particella che si tro-va nel punto P?
(Gara di 1° livello edizione 2004) Il seguente diagramma mostra un’onda periodica.
A B C F D E G 7 t s t t t s t a d e b c s s s x P s 6 Generatore di onde 4 m Muro 5 CAPITOLO 12 Le onde e il suono 365 esercizi Quale delle seguenti rappresenta una coppia di punti in fase?
a A e C b B e D c C ed F d E e G e A ed F (Gara di 1° livello edizione 2004) Il diagramma mostra una corda elastica tesa fra i punti fissi P e Q. Sulla corda è presente un’onda stazionaria.
Quale o quali delle seguenti affermazioni, riguardo ai due punti indicati X e Y della corda, sono corrette?
1) I due punti oscillano con una differenza di fase ugua-le a π.
2) I due punti hanno lo stesso periodo di oscillazione.
3) La distanza tra i due punti è uguale a una lunghezza d’onda.
a Tutte e tre.
b Solo la 1 e la 2.
c Solo la 2 e la 3.
d Solo la 1.
e Solo la 3.
(Gara di 1° livello edizione 2002) La figura mostra due impulsi,ciascuno di lunghezza l,che si muovono lungo una corda l’uno verso l’altro alla stes-sa velocità. Quale disegno rappresenta meglio la forma della corda quando entrambi raggiungono il tratto AB?
(Gara di 1° livello edizione 2007) e 15 mm d c 10 mm 20 mm 20 mm b a 10 mm 20 mm (Impulsi non in scala) 10 mm A B l l l 9 P X Q Y 8 TEST DI AMMISSIONE ALL’UNIVERSITÀ Due corde dello stesso materiale e con diverso diame-tro, una grossa e una sottile, sono collegate tra loro.
Un’onda viaggia sulla corda grossa e raggiunge la con-nessione con la corda sottile. Quale delle seguenti gran-dezze cambia alla connessione tra le due corde?
a Frequenza.
b Periodo.
c Nessuna.
d Velocità di propagazione.
(Concorso a borse di studio per l’iscrizione ai corsi di laurea della classe «Scienze e Tecnologie Fisiche» della SIF, 2006-2007) Il suono è legato all’emissione, propagazione, ricezione e percezione di onde. Quale delle seguenti affermazioni è corretta?
a Si tratta di onde elettromagnetiche.
b Se la frequenza è di 100 000 Hz, l’orecchio umano percepisce l’onda come suono.
c Se la frequenza è di 1 Hz, l’orecchio umano percepi-sce l’onda come suono.
d Se la frequenza è di 3000 Hz, l’orecchio umano per-cepisce l’onda come suono.
e Il suono si propaga in qualunque mezzo, compreso il vuoto assoluto.
(Prova di ammissione al corso di laurea in Odontoiatria e Protesi dentaria, 1999-2000) Il suono è un’onda che si propaga: a nel vuoto con velocità di 340 m/s.
b nel vuoto con frequenza uguale a 20 Hz.
c in un mezzo elastico con velocità che dipende dal mezzo.
d nel vuoto con velocità di 3 108 m/s.
e in un mezzo elastico con velocità uguale a 3 108 m/s.
(Prova di ammissione al corso di laurea in Architettura e Ingegneria Edile, 2000-2001) PROVE D’ESAME ALL’UNIVERSITÀ Un’onda sinusoidale che si propaga lungo una corda è descritta dall’equazione x(t) 0,03 sen (3,4 t). Si dedu-ce che la frequenza dell’onda è: a 0,54 Hz b 1,08 Hz c 5,4 Hz d 0,11 Hz (Esame di Fisica, corso di laurea in Infermieristica, Uni-versità di Napoli, 2005-2006) 1 3 2 1 CAPITOLO 12 Le onde e il suono 366 |
13741 | https://teachy.ai/es/resumenes/educacion-secundaria-es/secundaria-3-es/ciencias-es/unidad-astronomica-or-resumen-tradicional-90c | Resumen de Unidad Astronómica | Resumen Tradicional
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Unidad Astronómica
Resumen Tradisional | Unidad Astronómica
Contextualización
Al observar el cielo y estudiar el universo, los astrónomos se enfrentan a distancias tan enormes que medir en kilómetros resulta poco práctico. Para abordar este reto, los científicos recurren a unidades especiales que hacen más accesible la comprensión y comunicación de estas distancias. Una de estas unidades es la Unidad Astronómica (UA), que representa la distancia media entre la Tierra y el Sol, que equivale a aproximadamente 150 millones de kilómetros.
La Unidad Astronómica es una herramienta fundamental en astronomía, ya que permite realizar mediciones y comunicarse acerca de las distancias entre los cuerpos celestes del Sistema Solar de forma más clara. Sin esta unidad estándar, sería extremadamente complejo manejar cifras tan elevadas y entender las vastas dimensiones que se presentan. La UA ayuda a simplificar los cálculos y brinda una forma estandarizada de representar estas enormes distancias.
¡Para Recordar!
Definición de la Unidad Astronómica (UA)
La Unidad Astronómica (UA) es la distancia media entre la Tierra y el Sol, que equivale aproximadamente a 150 millones de kilómetros. Esta definición se estableció para simplificar la medición de distancias dentro del Sistema Solar, ya que medir en kilómetros puede dar lugar a números extremadamente grandes y difíciles de gestionar.
La UA hace que la comunicación y la comprensión de las distancias entre los cuerpos celestes sean más prácticas y eficientes. Sin esta unidad, realizar cálculos e interpretar datos sobre las posiciones y movimientos de los planetas sería un verdadero desafío.
Además, la UA es una unidad clave en astronomía, ya que sirve de base para otras mediciones de distancia, como el año luz y el parsec. Proporciona un punto de referencia común, crucial para la comparación y conversión entre diferentes unidades astronómicas.
La UA equivale a aproximadamente 150 millones de kilómetros.
Facilita la medición y comunicación de distancias en el Sistema Solar.
Sirve de base para otras unidades de medida astronómica.
Historia y Necesidad de la Unidad Astronómica
La necesidad de una unidad estándar como la Unidad Astronómica surgió debido a la inmensidad del espacio y la complejidad de trabajar con números muy elevados. Antes de la introducción de la UA, las mediciones en kilómetros eran poco prácticas y complicadas, lo que dificultaba el trabajo de los astrónomos y científicos.
Históricamente, la definición de la UA se basó en la órbita de la Tierra alrededor del Sol. Con el paso del tiempo, los avances tecnológicos y científicos han permitido realizar mediciones más precisas, consolidando la UA como una unidad esencial en la astronomía contemporánea.
La UA no solo simplifica la comunicación entre científicos, sino que también hace más fácil la enseñanza y comprensión de las escalas astronómicas para los estudiantes y el público. Es una herramienta indispensable que contribuye a la precisión y claridad en la investigación y estudios astronómicos.
Surgió de la necesidad de tratar distancias muy grandes.
Basada en la órbita de la Tierra alrededor del Sol.
Facilita la comunicación y comprensión de las escalas astronómicas.
Conversión de Distancias a Unidades Astronómicas
Convertir distancias de kilómetros a Unidades Astronómicas es un proceso sencillo que implica dividir la distancia en kilómetros por la distancia media de la Tierra al Sol (150 millones de kilómetros). Por ejemplo, para convertir 480 millones de kilómetros a UA, se divide 480 millones entre 150 millones, resultando aproximadamente en 3.2 UA.
Esta conversión es fundamental para facilitar la comprensión y la comunicación de las distancias dentro del Sistema Solar. Al emplear la UA, científicos y estudiantes pueden trabajar con cifras más manejables y comparables, lo que simplifica el análisis y el estudio de las distancias astronómicas.
Practicar la conversión de distancias a unidades astronómicas es una habilidad importante para los estudiantes, ya que refuerza su comprensión sobre la escala de las distancias en el espacio y fomenta el desarrollo de habilidades matemáticas aplicadas a la astronomía.
Divide la distancia en kilómetros por 150 millones para convertir a UA.
Facilita la comprensión y comunicación de las distancias en el Sistema Solar.
Promueve la comprensión de la escala de las distancias en el espacio.
Aplicaciones de la Unidad Astronómica
La Unidad Astronómica se utiliza ampliamente en astronomía para medir y comunicar distancias dentro del Sistema Solar. Por ejemplo, la distancia media de Marte al Sol es aproximadamente 1.5 UA, mientras que la de Júpiter es de alrededor de 5.2 UA. Estas mediciones ayudan a comprender mejor las posiciones y órbitas de los planetas.
Además de las distancias planetarias, la UA se aplica en diversas investigaciones astronómicas, como la determinación de las órbitas de asteroides y cometas, y en las misiones de sondas espaciales. Conocer las distancias en UA facilita la planificación de trayectorias y la realización de cálculos precisos para los viajes espaciales.
La UA también sirve como base de referencia para unidades de medida más grandes, como el año luz (la distancia que la luz recorre en un año) y el parsec (aproximadamente 3.26 años luz). Esta jerarquía de unidades permite una comprensión más extensa de las escalas astronómicas más allá del Sistema Solar.
Se utiliza para medir distancias en el Sistema Solar.
Facilita la determinación de órbitas y trayectorias de sondas espaciales.
Base de referencia para unidades más grandes como el año luz y el parsec.
Términos Clave
Unidad Astronómica: Distancia media entre la Tierra y el Sol, aproximadamente 150 millones de kilómetros.
Año luz: Distancia que la luz recorre en un año, aproximadamente 9.46 billones de kilómetros.
Parsec: Unidad de medida de distancia equivalente a aproximadamente 3.26 años luz.
Distancias astronómicas: Mediciones utilizadas para describir las distancias entre cuerpos celestes.
Conversión de distancia: Proceso de transformación de mediciones de kilómetros a Unidades Astronómicas.
Conclusiones Importantes
La Unidad Astronómica (UA) es una herramienta fundamental en astronomía que simplifica la medición y comunicación de vastas distancias en el espacio. Definida como la distancia media entre la Tierra y el Sol, aproximadamente 150 millones de kilómetros, la UA permite a científicos y estudiantes trabajar con números más manejables y comparables, facilitando el análisis y estudio de las escalas astronómicas.
La necesidad de la UA surgió por la inmensidad del espacio y la dificultad de trabajar con cifras extremadamente grandes utilizando kilómetros. La UA no solo simplifica la comunicación entre científicos, sino que también contribuye a la enseñanza y comprensión de las escalas astronómicas. Es una unidad clave que sirve de base para otras medidas, como el año luz y el parsec.
Las aplicaciones prácticas de la UA incluyen medir distancias entre planetas en el Sistema Solar y determinar las órbitas de asteroides y cometas. Comprender las distancias en UA es crucial para la planificación de trayectorias y cálculos precisos en misiones espaciales. Conocer la UA contribuye a explorar mejor el cosmos y al avance de la astronomía y campos relacionados.
Consejos de Estudio
Repasa los conceptos tratados en clase, centrándote en la definición y la importancia de la Unidad Astronómica.
Practica la conversión de distancias de kilómetros a Unidades Astronómicas utilizando ejemplos prácticos.
Lee artículos y libros sobre astronomía que aborden la aplicación de la Unidad Astronómica en diferentes contextos para profundizar tu conocimiento.
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13742 | https://harabor.net/data/papers/harabor-grastien-aaai11.pdf | Online Graph Pruning for Pathfinding on Grid Maps Daniel Harabor and Alban Grastien NICTA and The Australian National University Email: firstname.lastname@nicta.com.au Abstract Pathfinding in uniform-cost grid environments is a problem commonly found in application areas such as robotics and video games.
The state-of-the-art is dominated by hierar-chical pathfinding algorithms which are fast and have small memory overheads but usually return suboptimal paths. In this paper we present a novel search strategy, specific to grids, which is fast, optimal and requires no memory overhead. Our algorithm can be described as a macro operator which iden-tifies and selectively expands only certain nodes in a grid map which we call jump points. Intermediate nodes on a path connecting two jump points are never expanded. We prove that this approach always computes optimal solutions and then undertake a thorough empirical analysis, comparing our method with related works from the literature. We find that searching with jump points can speed up A by an order of magnitude and more and report significant improvement over the current state of the art.
Introduction Widely employed in areas such as robotics (Lee and Yu 2009), artificial intelligence (Wang and Botea 2009) and video games (Davis 2000; Sturtevant 2007), the ubiqui-tous undirected uniform-cost grid map is a highly popular method for representing pathfinding environments. Regu-lar in nature, this domain typically features a high degree of path symmetry (Harabor and Botea 2010; Pochter et al.
2010). Symmetry in this case manifests itself as paths (or path segments) which share the same start and end point, have the same length and are otherwise identical save for the order in which moves occur. Unless handled properly, sym-metry can force search algorithms to evaluate many equiva-lent states and prevents real progress toward the goal.
In this paper we deal with such path symmetries by devel-oping a macro operator that selectively expands only certain nodes from the grid, which we call jump points. Moving from one jump point to the next involves travelling in a fixed direction while repeatedly applying a set of simple neigh-bour pruning rules until either a dead-end or a jump point is reached. Because we do not expand any intermediate nodes between jump points our strategy can have a dramatic pos-itive effect on search performance. Furthermore, computed Copyright c ⃝2011, Association for the Advancement of Artificial Intelligence (www.aaai.org). All rights reserved.
solutions are guaranteed to be optimal. Jump point pruning is fast, requires no preprocessing and introduces no mem-ory overheads. It is also largely orthogonal to many existing speedup techniques applicable to grid maps.
We make the following contributions: (i) a detailed de-scription of the jump points algorithm; (ii) a theoretical re-sult which shows that searching with jump points preserves optimality; (iii) an empirical analysis comparing our method with two state-of-the-art search space reduction algorithms.
We run experiments on a range of synthetic and real-world benchmarks from the literature and find that jump points improve the search time performance of standard A by an order of magnitude and more.
We also report signif-icant improvement over Swamps (Pochter et al. 2010), a recent optimality preserving pruning technique, and perfor-mance that is competitive with, and in many cases domi-nates, HPA (Botea, M¨ uller, and Schaeffer 2004); a well known sub-optimal pathfinding algorithm.
Related Work Approaches for identifying and eliminating search-space symmetry have been proposed in areas including planning (Fox and Long 1999), constraint programming (Gent and Smith 2000), and combinatorial optimization (Fukunaga 2008). Very few works however explicitly identify and deal with symmetry in pathfinding domains such as grid maps.
Empty Rectangular Rooms (Harabor and Botea 2010) is an offline symmetry breaking technique which attempts to redress this oversight. It decomposes grid maps into a se-ries of obstacle-free rectangles and replaces all nodes from the interior of each rectangle with a set of macro edges that facilitate optimal travel. Specific to 4-connected maps, this approach is less general than jump point pruning. It also re-quires offline pre-processing whereas our method is online.
The dead-end heuristic (Bj¨ ornsson and Halld´ orsson 2006) and Swamps (Pochter et al. 2010) are two similar pruning techniques related to our work. Both decompose grid maps into a series of adjacent areas. Later, this decomposition is used to identify areas not relevant to optimally solving a particular pathfinding instance. This objective is similar yet orthogonal to our work where the aim is to reduce the effort required to explore any given area in the search space.
A different method for pruning the search space is to identify dead and redundant cells (Sturtevant, Bulitko, and Bj¨ ornsson 2010).
Developed in the context of learning-based heuristic search, this method speeds up search only after running multiple iterations of an iterative deepening al-gorithm. Further, the identification of redundant cells re-quires additional memory overheads which jump points do not have.
Fast expansion (Sun et al. 2009) is another related work that speeds up optimal A search. It avoids unnecessary open list operations when it finds a successor node just as good (or better) than the best node in the open list. Jump points are a similar yet fundamentally different idea: they allow us to identify large sets of nodes that would be ordi-narily expanded but which can be skipped entirely.
In cases where optimality is not required, hierarchical pathfinding methods are pervasive. They improve perfor-mance by decomposing the search space, usually offline, into a much smaller approximation. Algorithms of this type, such as HPA (Botea, M¨ uller, and Schaeffer 2004), are fast and memory-efficient but also suboptimal.
Notation and Terminology We work with undirected uniform-cost grid maps.
Each node has ≤8 neighbours and is either traversable or not.
Each straight (i.e.
horizontal or vertical) move, from a traversable node to one of its neighbours, has a cost of 1; diagonal moves cost √ 2. Moves involving non-traversable (obstacle) nodes are disallowed. The notation ⃗ d refers to one of the eight allowable movement directions (up, down, left, right etc.). We write y = x+k⃗ d when node y can be reached by taking k unit moves from node x in direction ⃗ d. When ⃗ d is a diagonal move, we denote the two straight moves at 45 deg to ⃗ d as ⃗ d1 and ⃗ d2.
A path π = ⟨n0, n1, . . . , nk⟩is a cycle-free ordered walk starting at node n0 and ending at nk. We will sometimes use the setminus operator in the context of a path: for example π \x. This means that the subtracted node x does not appear on (i.e. is not mentioned by) the path. We will also use the function len to refer the length (or cost) of a path and the function dist to refer to the distance between two nodes on the grid: e.g. len(π) or dist(n0, nk) respectively.
Jump Points In this section we introduce a search strategy for speeding up optimal search by selectively expanding only certain nodes on a grid map which we term jump points. We give an ex-ample of the basic idea in Figure 1(a).
Here the search is expanding a node x which has as its par-ent p(x); the direction of travel from p(x) to x is a straight move to the right. When expanding x we may notice that there is little point to evaluating any neighbour highlighted grey as the path induced by such a move is always domi-nated by (i.e. no better than) an alternative path which men-tions p(x) but not x. We will make this idea more precise in the next section but for now it is sufficient to observe that the only non-dominated neighbour of x lies immedi-ately to the right. Rather than generating this neighbour and adding it to the open list, as in the classical A algorithm, Figure 1: Examples of straight (a) and diagonal (b) jump points. Dashed lines indicate a sequence of interim node evaluations that reached a dead end. Strong lines indicate eventual successor nodes.
Figure 2: We show several cases where a node x is reached from its parent p(x) by either a straight or diagonal move.
When x is expanded we can prune from consideration all nodes marked grey.
we propose to simply step to the right and continue mov-ing in this direction until we encounter a node such as y; which has at least one other non-dominated neighbour (here z). If we find a node such as y (a jump point) we generate it as a successor of x and assign it a g-value (or cost-so-far) of g(y) = g(x) + dist(x, y). Alternatively, if we reach an obstacle we conclude that further search in this direction is fruitless and generate nothing.
In the remainder of this section we will develop a macro-step operator which speeds up node expansion by identify-ing jump point successors in the case of both straight and diagonal moves. First it will be necessary to define a series of pruning rules to determine whether a node should be gen-erated or skipped. This will allow us to make precise the notion of a jump point and give a detailed description of the jump points algorithm. Then, we prove that the process of “jumping” nodes, such as x in Figure 1(a), has no effect on the optimality of search.
Neighbour Pruning Rules In this section we develop rules for pruning the set of nodes immediately adjacent to some node x from the grid. The ob-jective is to identify from each set of such neighbours, i.e.
neighbours(x), any nodes n that do not need to be evalu-ated in order to reach the goal optimally. We achieve this by comparing the length of two paths: π, which begins with node p(x) visits x and ends with n and another path π′ which also begins at node p(x) and ends with n but does not men-tion x. Additionally, each node mentioned by either π or π′ must belong to neighbours(x).
There are two cases to consider, depending on whether the transition to x from its parent p(x) involves a straight move or a diagonal move. Note that if x is the start node p(x) is null and nothing is pruned.
Straight Moves: We prune any node n ∈neighbours(x) which satisfies the following dominance constraint: len( ⟨p(x), . . . , n⟩\ x ) ≤len( ⟨p(x), x, n⟩) (1) Figure 2(a) shows an example. Here p(x) = 4 and we prune all neighbours except n = 5.
Diagonal Moves: This case is similar to the pruning rules we developed for straight moves; the only difference is that the path which excludes x must be strictly dominant: len( ⟨p(x), . . . , n⟩\ x ) < len( ⟨p(x), x, n⟩) (2) Figure 2(c) shows an example. Here p(x) = 6 and we prune all neighbours except n = 2, n = 3 and n = 5.
Assuming neighbours(x) contains no obstacles, we will refer to the nodes that remain after the application of straight or diagonal pruning (as appropriate) as the natural neigh-bours of x. These correspond to the non-gray nodes in Fig-ures 2(a) and 2(c). When neighbours(x) contains an obsta-cle, we may not be able to prune all non-natural neighbours.
If this occurs we say that the evaluation of each such neigh-bour is forced.
Definition 1. A node n ∈neighbours(x) is forced if: 1. n is not a natural neighbour of x 2. len( ⟨p(x), x, n⟩) < len( ⟨p(x), . . . , n⟩\ x ) In Figure 2(b) we show an example of a straight move where the evaluation of n = 3 is forced. Figure 2(d) shows an sim-ilar example involving a diagonal move; here the evaluation of n = 1 is forced.
Algorithmic Description We begin by making precise the concept of a jump point.
Definition 2. Node y is the jump point from node x, heading in direction ⃗ d, if y minimizes the value k such that y = x+k⃗ d and one of the following conditions holds: 1. Node y is the goal node.
2. Node y has at least one neighbour whose evaluation is forced according to Definition 1.
3. ⃗ d is a diagonal move and there exists a node z = y +ki⃗ di which lies ki ∈N steps in direction ⃗ di ∈{ ⃗ d1, ⃗ d2} such that z is a jump point from y by condition 1 or condition 2.
Figure 1(b) shows an example of a jump point which is identified by way of condition 3. Here we start at x and travel diagonally until encountering node y. From y, node z Algorithm 1 Identify Successors Require: x: current node, s: start, g: goal 1: successors(x) ←∅ 2: neighbours(x) ←prune(x, neighbours(x)) 3: for all n ∈neighbours(x) do 4: n ←jump(x, direction(x, n), s, g) 5: add n to successors(x) 6: return successors(x) can be reached with ki = 2 horizontal moves. Thus z is a jump point successor of y (by condition 2) and this in turn identifies y as a jump point successor of x The process by which individual jump point successors are identified is given in Algorithm 1. We start with the pruned set of neighbours immediately adjacent to the cur-rent node x (line 2). Then, instead of adding each neighbour n to the set of successors for x, we try to “jump” to a node that is further away but which lies in the same relative di-rection to x as n (lines 3:5). For example, if the edge (x, n) constitutes a straight move travelling right from x, we look for a jump point among the nodes immediately to the right of x. If we find such a node, we add it to the set of succes-sors instead of n. If we fail to find a jump point, we add nothing. The process continues until the set of neighbours is exhausted and we return the set of successors for x (line 6).
In order to identify individual jump point successors we will apply Algorithm 2. It requires an initial node x, a direc-tion of travel ⃗ d, and the identities of the start node s and the goal node g. In rough overview, the algorithm attempts to establish whether x has any jump point successors by step-ping in the direction ⃗ d (line 1) and testing if the node n at that location satisfies Definition 2. When this is the case, n is designated a jump point and returned (lines 5, 7 and 11).
When n is not a jump point the algorithm recurses and steps again in direction ⃗ d but this time n is the new initial node (line 12). The recursion terminates when an obstacle is en-countered and no further steps can be taken (line 3). Note that before each diagonal step the algorithm must first fail to detect any straight jump points (lines 9:11). This check corresponds to the third condition of Definition 2 and is es-sential for preserving optimality.
Optimality In this section we prove that for each optimal length path in a grid map there exists an equivalent length path which can be found by only expanding jump point nodes during search (Theorem 1). Our result is derived by identifying for each optimal path a symmetric alternative which we split into contiguous segments. We then prove that each turning point along this path is also a jump point.
Algorithm 2 Function jump Require: x: initial node, ⃗ d: direction, s: start, g: goal 1: n ←step(x, ⃗ d) 2: if n is an obstacle or is outside the grid then 3: return null 4: if n = g then 5: return n 6: if ∃n′ ∈neighbours(n) s.t. n′ is forced then 7: return n 8: if ⃗ d is diagonal then 9: for all i ∈{1, 2} do 10: if jump(n, ⃗ di, s, g) is not null then 11: return n 12: return jump(n, ⃗ d, s, g) Definition 3. A turning point is any node ni along a path where the direction of travel from the previous node ni−1 to ni is different to the direction of travel from ni to the subsequent node ni+1.
Figure 3 depicts the three possible kinds of turning points which we may encounter on an optimal path. A diagonal-to-diagonal turning point at node nk (Figure 3(a)) involves a diagonal step from its parent nk−1 followed by a second diagonal step, this time in a different direction, from nk to its successor nk+1. Similarly, a straight-to-diagonal (or diagonal-to-straight) turning point involves a straight (diag-onal) step from nk−1 to reach nk followed by a diagonal (straight) step to reach its successor nk+1 (Figure 3(b) and 3(c) respectively). Other types of turning points, such as straight-to-straight, are trivially suboptimal and not consid-ered here (they are pruned by the rules we developed earlier; see again Figure 2).
We are now ready to develop an equivalence relation between jump points and the turning points that appear along certain optimal length symmetric paths which we term diagonal-first.
Definition 4. A path π is diagonal-first if it contains no straight-to-diagonal turning point ⟨nk−1, nk, nk+1⟩which could be replaced by a diagonal-to-straight turning point ⟨nk−1, n′ k, nk+1⟩s.t. the length of π remains unchanged.
Given an arbitrary optimal length path π, we can always derive a symmetric diagonal-first path π′ by applying Algo-rithm 3 to π. Note that this is merely as a conceptual device.
Lemma 1. Each turning point along an optimal diagonal-first path π′ is also a jump point.
Proof. Let nk be an arbitrary turning point node along π′.
We will consider three cases, each one corresponding to one of the three possible kinds of optimal turning points illus-trated in Figure 3.
Diagonal-to-Diagonal: Since π′ is optimal, there must be an obstacle adjacent to both nk and nk−1 which is forc-ing a detour. We know this because if there were no obsta-cle we would have dist(nk−1, nk+1) < dist(nk−1, nk) + dist(nk, nk+1) which contradicts the fact that π′ is optimal.
We conclude that nk+1 is a forced neighbour of nk. This is sufficient to satisfy the second condition of Definition 1, making nk a jump point.
Algorithm 3 Compute Diagonal First Path Require: π: an arbitrary optimal length path 1: select an adjacent pair of edges appearing along π where (nk−1, nk) is a straight move and (nk, nk+1) is a diag-onal move.
2: replace (nn−1, nk) and (nk, nk+1) with two new edges: (nk−1, n′ k), which is a diagonal move and (n′ k, nk+1) which is a straight move. The operation is successful if (nk−1, n′ k) and (n′ k, nk+1) are both valid moves; i.e.
node n′ k is not an obstacle.
3: repeat lines 1 and 2, selecting and replacing adjacent edges, until no further changes can be made to π.
4: return π Figure 3: Types of optimal turning points. (a) Diagonal-to-Diagonal (b) Straight-to-Diagonal (c) Diagonal-to-Straight.
Straight-to-Diagonal: In this case there must be an ob-stacle adjacent to nk. If this were not true nk could be re-placed by a Diagonal-to-Straight turning point which contra-dicts the fact that π′ is diagonal-first. Since π′ is guaranteed to be diagonal-first we derive the fact that nk+1 is a forced neighbour of nk. This satisfies the second condition of Def-inition 1 and we conclude nk is a jump point.
Diagonal-to-Straight: There are two possibilities in this case, depending on whether the goal is reachable by a series of straight steps from nk or whether π′ has additional turning points. If the goal is reachable by straight steps then nk has a jump point successor which satisfies the third condition of Definition 1 and we conclude nk is also a jump point. If nk is followed by another turning point, nl, then that turning point must be Straight-to-Diagonal and, by the argument for that case, also a jump point. We again conclude that nk has a jump point successor which satisfies the third condition of Definition 1. Thus, nk is also a jump point.
Theorem 1. Searching with jump point pruning always re-turns an optimal solution.
Proof. Let π be an arbitrarily chosen optimal path between two nodes on a grid and π′ a diagonal-first symmetric equiv-alent which is derived by applying Algorithm 3 to π. We will show that every turning point mentioned by π′ is ex-panded optimally when searching with jump point pruning.
We argue as follows: Divide π′ into a series of adjacent segments s.t. π′ = π′ 0 + π′ 1 + . . . + π′ n. Each π′ i = ⟨n0, n1, . . . , nk−1, nk⟩is a subpath along which all moves involve travelling in the same direction (e.g. only “up” or “down” etc). Notice that with the exception of the start and goal, every node at the beginning and end of a segment is also a turning point.
Since each π′ i consists only of moves in a single direc-tion (straight or diagonal) we can use Algorithm 2 to jump from n0 ∈π′ i, the node at beginning of each segment to nk ∈π′ i, the node at the end, without necessarily stopping to expand every node in between. Intermediate expansions may occur but the fact that we reach nk optimally from n0 is guaranteed. It remains to show only that both n0 and nk are identified as jump points and thus necessarily expanded.
By Lemma 1 each turning point along π′ is also a jump point, so every turning point node must be expanded dur-ing search. Only the start and goal remain. The start node is necessarily expanded at the beginning of each search while the goal node is a jump point by definition. Thus both are expanded.
Experimental Setup We evaluate the performance of jump point pruning on four benchmarks taken from the freely available pathfinding library Hierarchical Open Graph (HOG,
googlecode.com/p/hog2): • Adaptive Depth is a set of 12 maps of size 100×100 in which approximately 1 3 of each map is divided into rect-angular rooms of varying size and a large open area inter-spersed with large randomly placed obstacles. For this benchmark we randomly generated 100 valid problems per map for a total of 1200 instances.
• Baldur’s Gate is a set of 120 maps taken from BioWare’s popular roleplaying game Baldur’s Gate II: Shadows of Amn; it appears regularly as a standard benchmark in the literature (Bj¨ ornsson and Halld´ orsson 2006; Harabor and Botea 2010; Pochter et al.
2010).
We use the varia-ton due to Nathan Sturtevant where all maps have been scaled to size 512×512 to more accurately represent mod-ern pathfinding environments. Maps and all instances are available from • Dragon Age is another realistic benchmark; this time taken from BioWare’s recent roleplaying game Dragon Age: Origins. It consists of 156 maps ranging in size from 30 × 21 to 1104 × 1260. For this benchmark we used a large set of randomly generated instances, again due to Nathan Sturtevant and available from http:// movingai.com.
• Rooms is a set of 300 maps of size 256×256 which are divided into symmetric rows of small rectangular ar-eas (7 × 7), connected by randomly placed entrances.
This benchmark has previously appeared in (Pochter et al.
2010). For this benchmark we randomly generated 100 valid problems per map for a total of 30000 instances.
Our test machine is a 2.93GHz Intel Core 2 Duo processor with 4GB RAM running OSX 10.6.4.
Results To evaluate jump points we use a generic implementation of A which we adapted to facilitate online neighbour prun-ing and jump point identification. We discuss performance in terms of speedup: i.e. relative improvement to the time taken to solve a given problem, and the number of nodes expanded, when searching with and without graph pruning applied to the grid. Using this metric a search time speedup of 2.0 is twice as fast while a node expansion speedup of 2.0 indicates half the number of nodes were expanded. In each case higher is better. Figure 4 shows average search time speedups across each of our four benchmarks. Table 1 shows average node expansion speedups; the best results for each column are in bold.
Comparison with Swamps: We begin by comparing jump points to Swamps (Pochter et al. 2010): an optimal-ity preserving pruning technique for speeding up pathfind-ing. We used the authors’ source code, and their imple-mentation of A, and ran all experiments using their rec-ommended running parameters: a swamp seed radius of 6 and “no change limit” of 2.
A. Depth B. Gate D.Age Rooms Jump Points 20.37 215.36 35.95 13.41 Swamps 1.89 2.44 2.99 4.70 HPA 4.14 9.37 9.63 5.11 Table 1: Average A node expansion speedup.
As per Figure 4, jump point pruning shows a convinc-ing improvement to average search times across all bench-marks.
The largest differences are observed on Baldur’s Gate and Dragon Age where searching with jump points reaches the goal between 25-30 times sooner while search-ing with Swamps attains only a 3-5 times speedup. Similar trends are observed when looking at Table 1, where the im-provement to the total number of nodes expanded is even more pronounced.
Based on these results we conclude that while Swamps are effective for identifying areas of the map not relevant to reaching the goal, those areas which remain still require sig-nificant effort to search. Jump points use a much stronger yet orthogonal strategy to prove that many nodes expanded by Swamps can be ignored. As the two ideas appear comple-mentary we posit that they could be easily combined: first, apply a Swamps-based decomposition to prune areas not rel-evant to the current search. Then, use jump points to search the remaining portions of the map.
Comparison with HPA: Next, we compare jump points pruning to the HPA algorithm (Botea, M¨ uller, and Schaeffer 2004). Though sub-optimal, HPA is very fast and widely applied in video games. To evaluate HPA we measure the total cost of insertion and hierarchical search.
We did not refine any abstract paths, assuming instead that a database of pre-computed intra-cluster paths was available.
Such a configuration represents the fastest generic imple-mentation of HPA but requires additional memory over-heads. While searching we used a single-level abstraction hierarchy and a fixed cluster size of 10. These settings are recommended by the original authors who note that larger clusters and more levels are often of little benefit.
As per Figure 4, jump points pruning is shown to be highly competitive with HPA. On Adaptive Depth, Bal-dur’s Gate and Rooms, jump points have a clear advantage – improving on HPA search times by several factors in some cases. On the Dragon Age benchmark jump points have a small advantage for problems of length < 500 but for longer instances there is very little between the two algorithms. Ta-ble 1 provides further insight: although searching with jump points expands significantly fewer nodes than HPA, each such operation takes longer.
We conclude that jump point pruning is a competitive sub-stitute for HPA and, for a wide variety of problems, can help to find solutions significantly faster. HPA can still be advantageous, particularly if a memory overhead is accept-able and optimality is not important, but only if the cost of inserting the start and goal nodes into the abstract graph (and possibly refinement, if a path database is not available), can be sufficiently amortized over the time required to find an abstract path. One direction for further work could be to use 50 100 150 0 5 10 15 20 25 30 35 Search Time Speedup: Adaptive Depth Path Length Average Speedup Factor Jump Points Swamps HPA 0 100 200 300 400 500 0 5 10 15 20 25 30 35 Search Time Speedup: Baldur's Gate (512x512) Path Length Average Speedup Factor Jump Points Swamps HPA 0 500 1000 1500 2000 2500 0 5 10 15 20 25 30 35 Search Time Speedup: Dragon Age Path Length Average Speedup Factor Jump Points Swamps HPA 0 100 200 300 400 500 0 5 10 15 20 25 30 35 Search Time Speedup: Rooms Path Length Average Speedup Factor Jump Points Swamps HPA Figure 4: Average A search time speedup on our each of our four benchmarks.
jump point pruning to speed up HPA: for example during insertion and refinement.
Conclusion We introduce a new online node pruning strategy for speed-ing up pathfinding on undirected uniform-cost grid maps.
Our algorithm identifies and selectively expands only certain nodes from a grid map which we call jump points. Moving between jump points involves only travelling in a fixed di-rection, either straight or diagonal. We prove that interme-diate nodes on a path between two jump points never need to be expanded and “jumping” over them does not affect the optimality of search.
Our method is unique in the pathfinding literature in that it has very few disadvantages: it is simple, yet highly effective; it preserves optimality, yet requires no extra memory; it is extremely fast, yet requires no preprocessing. Further, it is largely orthogonal to and easily combined with competing speedup techniques from the literature. We are unaware of any other algorithm which has all these features.
The new algorithm is highly competitive with re-lated works from the literature.
When compared to Swamps (Pochter et al. 2010), a recent state-of-the-art op-timality preserving pruning technique, we find that jump points are up to an order of magnitude faster.
We also show that jump point pruning is competitive with, and in many instances clearly faster than, HPA; a popular though sub-optimal pathfinding technique often employed in perfor-mance sensitive applications such as video games.
One interesting direction for further work is to extend jump points to other types of grids, such as hexagons or texes (Yap 2002). We propose to achieve this by develop-ing a series of pruning rules analogous to those given for square grids. As the branching factor on these domains is lower than square grids, we posit that jump points could be even more effective than observed in the current paper.
Another interesting direction is combining jump points with other speedup techniques: e.g. Swamps or HPA.
Acknowledgements We would like to thank Adi Botea, Philip Kilby and Patrik Haslum for their encouragement, support and many helpful suggestions during the development of this work. We also thank Nir Pochter for kindly providing us with source code for the Swamps algorithm.
NICTA is funded by the Australian Government as rep-resented by the Department of Broadband, Communications and the Digital Economy and the Australian Research Coun-cil through the ICT Centre of Excellence program.
References Bj¨ ornsson, Y., and Halld´ orsson, K. 2006. Improved heuris-tics for optimal path-finding on game maps. In AIIDE, 9–14.
Botea, A.; M¨ uller, M.; and Schaeffer, J. 2004. Near optimal hierarchical path-finding. J. Game Dev. 1(1):7–28.
Davis, I. L. 2000. Warp speed: Path planning for Star Trek Armada. In AAAI Spring Symposium (AIIDE), 18–21.
Fox, M., and Long, D. 1999. The detection and exploitation of symmetry in planning problems. In IJCAI, 956–961.
Fukunaga, A. S. 2008. Integrating symmetry, dominance, and bound-and-bound in a multiple knapsack solver.
In CPAIOR, 82–96.
Gent, I. P., and Smith, B. M. 2000. Symmetry breaking in constraint programming. In ECAI, 599–603.
Harabor, D., and Botea, A. 2010. Breaking path symmetries in 4-connected grid maps. In AIIDE, 33–38.
Lee, J.-Y., and Yu, W. 2009. A coarse-to-fine approach for fast path finding for mobile robots. In IROS, 5414 –5419.
Pochter, N.; Zohar, A.; Rosenschein, J. S.; and Felner, A.
2010. Search space reduction using swamp hierarchies. In AAAI.
Sturtevant, N. R.; Bulitko, V.; and Bj¨ ornsson, Y. 2010. On learning in agent-centered search. In AAMAS, 333–340.
Sturtevant, N. R. 2007. Memory-efficient abstractions for pathfinding. In AIIDE, 31–36.
Sun, X.; Yeoh, W.; Chen, P.-A.; and Koenig, S. 2009. Sim-ple optimization techniques for a-based search. In AAMAS (2), 931–936.
Wang, K.-H. C., and Botea, A. 2009. Tractable multi-agent path planning on grid maps. In IJCAI, 1870–1875.
Yap, P. 2002. Grid-based path-finding. In Canadian AI, 44–55. |
13743 | https://content.byui.edu/file/b8b83119-9acc-4a7b-bc84-efacf9043998/1/Math-1-5-7.html | Pie Charts
HomeMathWritingLanguageSkillsTutoring
Division and Percentages:
Pie Charts
In this lesson, you will practice interpreting pie charts. What is a pie?
A pie is a food dish that usually contains either fruit, meat, or vegetables and has a bottom and sometimes a top made out of pastry. It is usually round and divided into triangle-like slices, starting from the center to the outside edge. The picture above is of a pumpkin pie.
A pie chart is a round graph that uses pie-shaped wedges to represent the percentages of the whole.
Video Source (09:04 mins) | Transcript
Pie charts help us see how things compare to each other.
The bigger the slice of the pie chart, the higher the percentage of the whole it represents. The smaller the slice in the pie chart, the lower the percentage. The following pie chart represents the favorite colors of a group of people. The largest group of people represented by this pie chart like the color blue. Yellow is the least liked of the group.
The practice problems in this section help walk you through how to find percentages using data and a pie chart. You will need to know how to do this for the application activity.
Additional Information
Division:
If we want to convert a certain number of hours in a day to a percentage, we would divide that number by 24 hours. For example, if we wanted to know how much of the day is represented by 8.4 hours, we would divide 8.4 hours by 24 hours. The result is
8.4 hours 24 hours=0.35
Which we can write as 35%.
Notice that if we are dividing to calculate a percentage, the units in the numerator (hours) needs to be the same as the units in the denominator (hours). You can think of the units (hours) as “canceling” each other out.
Multiplication:
If we know that something is, say, 35% of a whole, then we can multiply 35% (or 0.35) by the total amount to find out how much 35% is of the total. For example, we learned that Daniel spends 35% of his day at work. This means that the number of hours he spent at work is:
0.35 × 24 hours = 8.4 hours.
Practice Problems
Answer the following questions.
The following chart shows how Daniel, a PathwayConnect student, uses his time in a typical weekday.
Everything Daniel does is represented in one of the slices in the chart, and the size of the slices indicates the percentage of time he does each activity. Larger slices indicate a greater amount of time.
Every day is comprised of 24 hours. This represents all (or 100%) of Daniel's time. So, the percentages of the slices on a pie chart must add up to 100%.
The word "percent" comes from Latin words that mean "for every hundred" or "one part in every hundred." You may find that you are already very familiar with percents. For examples, tithing represents 10% of our income, or 10 parts in every 100. So, a tithe (or 10%) of $200 would be $20.
Consider Daniel's pie chart as you complete the following questions.
When compared to the other activities, in what activity does Daniel spend the greatest portion of his time?
On what activity does Daniel spend the least amount of time?
What percentage of time does Daniel spend either sleeping or working?
What percentage of time is Daniel with family and friends?
How many hours does Daniel spend at work?
View Solutions
Solutions
Work (Written Solution) Daniel spends more time working than he does in any other activity. He works 35% of the time. 35% is greater than any other percentage illustrated on the graph.
Suppose we divided the day into 100 equal parts, Daniel would spend 35 of those working. If we were to consider all of Daniel’s time in a day, it would be 100% of his time. He spends 35% of his time working.
2. Shower and Dress (Written Solution) The smallest slice of the pie chart represents the percentage of time Daniel spends showering and dressing. The chart indicates that Daniel spent 1% of his day doing these things.
3. 60% (Solution Video | Transcript | Written Solution) The total percentage of time Daniel spends either sleeping or working can be calculated by adding the percentage of time he spends sleeping to the percentage of time he spends working. Daniel spends 25% of his time sleeping and 35% of his time working. Combining these, we get
25% + 35% = 60%
Notice that more than half of his time is spent either sleeping or working. Half of his time would be represented by 50%. Since 60% is greater than 50%, he spends more than half of his time either sleeping or working.
4. 9% (Written Solution) If we think about all the things Daniel can do in a day, all of this time must add up to 100%. When we think in terms of percentages, 100% represents the whole (or the entire amount.)
Daniel spends 25% of his time sleeping, 1% showering and dressing, 4% in scripture study and prayer, 5% travelling, 35% at work, 15% cooking and eating, and 6% doing homework. If we add these up, we find the total percentage of time he spends on these activities is:
25% + 1% + 4% + 5% + 35% + 15% + 6% = 91%.
The entire day is represented by 100%. So, if a full day is 100%, and 91% of the time is spent on other activities, the amount of time Daniel spends with family and friends is:
100% − 91% = 9%.
5. 8.4 hours (Solution Video | Transcript | Written Solution) In a full day, there are 24 hours. This represents 100% of the time Daniel has. The portion of the day that Daniel spends at work is 35%.
We want to know what portion of Daniel’s day is spent at work. In math, including with story problems, the word “of” can be a hint that we need to multiply. Notice what happens when we express the sentence “Daniel spends 35% of 24 hours at work.” In this case, the word “of” suggests multiplication.
We can convert this percentage to hours by multiplying 35% by 24 hours. Before we do this multiplication, it is helpful to rewrite 35% as 0.35.
We divide 35 by 100 to get 0.35. 35% and 0.35 are two ways to represent the same value.
Next, we multiply this by the total number of hours in a day:
0.35 × 24 hours = 8.4 hours
Daniel spends 8.4 hours each day at work. |
13744 | https://www.khanacademy.org/v/x-and-y-intercepts | Intercepts from an equation (video) | Khan Academy
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Algebra 1
Course: Algebra 1>Unit 4
Lesson 4: x-intercepts and y-intercepts
Intro to intercepts
x-intercept of a line
Intercepts from a graph
Intercepts from an equation
Intercepts from an equation
Intercepts from a table
Intercepts from a table
Intercepts of lines review (x-intercepts and y-intercepts)
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x-intercepts and y-intercepts
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Intercepts from an equation
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About About this video Transcript
Let's graph -5x + 4y = 20 from its intercepts. Intercepts are the places where a line crosses the x- and y-axes. When a line crosses the x-axis, the y value is 0. That's the x-intercept. When a line crosses the y-axis, the x value is 0. That's the y-intercept. The line passes through those points.Created by Sal Khan and Monterey Institute for Technology and Education.
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blanabas 7 years ago Posted 7 years ago. Direct link to blanabas's post “I'm a little confused I a...” more I'm a little confused I attempted this question from the next practice section:
Determine the intercepts of the line.
−4x+7=2y-3
And I came up with y= (0, -5) x= (2.5, 0)
and when I entered my answer it told me it was incorrect and I looked at the hint and it seems like they wanted the same coordinates just with positive numbers instead.
Is there a reason for that, for instance if I were taking the TASC do you think they'd accept my answer as well?
I remember reading there is basically an infinite amount of answers for Linear Equations, I'm just curious as to if I computed it on my own wrong or just an awkward way.
I did this:
-4x+7=2y-3
-7 -7
-4x = 2y -10
then
4(0) = 2y = 10
-10/2y = -5
y = -5
and
-10/-4 = 2.5
Sorry for the awful formatting my keyboard/ computer is all but completely broken
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TheLuckyRobot 7 years ago Posted 7 years ago. Direct link to TheLuckyRobot's post “-10/2y = -5 When you'r...” more -10/2y = -5
When you're dividing, you don't reverse the sign. It should be
2y - 10 = 0
+10 to both sides
2y = 10
divide both sides by positive 2 (because the sign doesn't change as we're dividing)
y = 10/2
and finally:
y = 5
Any other questions? Feel free to comment me back!
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Howard Mellow 5 years ago Posted 5 years ago. Direct link to Howard Mellow's post “My question is more from ...” more My question is more from experience. I am a 70 year old, taking this for a refresher. I know I have the right answer since I have attended a university in the field Electrical Engineering and Computer Science. However, I didn't get the problem right or wrong. The problem is how the computer accepts the answer. Do I have to enter the answer as a decimal, or a fraction? Please help since neither answer gave me the correct solutiopn. i.e. 3/4, 6/8, .75
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Pardiwala Khyaat 5 years ago Posted 5 years ago. Direct link to Pardiwala Khyaat's post “Hi Mr Mellow, I'm not sur...” more Hi Mr Mellow, I'm not sure of which questions you did, but certain questions would have certain formats for the answers. Some questions would want you to write the answer as a decimal, or another question would want you to type it as a fraction.
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ponyog 7 years ago Posted 7 years ago. Direct link to ponyog's post “What if you have a repeat...” more What if you have a repeating decimal in your intercept? What can you do to avoid it?
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Kim Seidel 7 years ago Posted 7 years ago. Direct link to Kim Seidel's post “An intercept can be any r...” more An intercept can be any real number: an integer, a fraction, a mixed number, a decimal.
In the case of a repeating decimal, I would keep the number as a fraction rather than changing to a decimal. It's an easy way to keep the entire value of the intercept.
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hernande.agustin647 5 years ago Posted 5 years ago. Direct link to hernande.agustin647's post “I am a 9th grader at Whit...” more I am a 9th grader at Whittier high school doing algebra 1, we are reviewing how to get Intercepts from an equation. One of the questions given is a little confusing for me and a few of my classmates. How would I get the X and Y intercepts from this? y−6=4(x+5)
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KC 5 years ago Posted 5 years ago. Direct link to KC's post “To get the X intercept pl...” more To get the X intercept plug in 0 for y so: 0-6=4(x+5)
Then solve for x the x intercept will be (-6.5, 0)
To get the Y intercept plug in 0 for x so: y-6=4(0+5)
then solve for y, the y interpect will be (0, 26)
I'm pretty sure, I hope this makes sense I can explain it in more depth if necessary.
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Anna 12 years ago Posted 12 years ago. Direct link to Anna's post “if the x intercept is 2 a...” more if the x intercept is 2 and the y intercept is -1 what is the slope of the corresponding line?
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1140858 4 years ago Posted 4 years ago. Direct link to 1140858's post “ more Slope's formula
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Lyze Of Kiel 🟢Read bio I'm alive 2 years ago Posted 2 years ago. Direct link to Lyze Of Kiel 🟢Read bio I'm alive's post “I can't seem to understan...” more I can't seem to understand this at all, the video doesn't seem to explain everything, can someone please explain this to me in layman's terms?
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David Severin 2 years ago Posted 2 years ago. Direct link to David Severin's post “An intercept is on the ax...” more An intercept is on the axis, so it will always have a 0 somewhere, the x intercept has y=0, so the point will be (#,0), and the y intercept has x=0, so the point will be (0,#). So given any equation, 3x + 4y = 24, we can find the x intercept by setting y=0, so we do 3x+4(0)=24 which gives 3x=24, then divide by 3 to get x=8. For y intercept, set x=0, 3(0)+4y=24, 4y=24, and divide by 4 to get y=6. Does this help?
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Benjamin Evans 2 years ago Posted 2 years ago. Direct link to Benjamin Evans's post “How do you know which poi...” more How do you know which point to plot first. In the last quis I had the right coordinated but in a different order.
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Kim Seidel 2 years ago Posted 2 years ago. Direct link to Kim Seidel's post “Ordered pairs are always ...” more Ordered pairs are always in the form (x-value, y-value). Make sure you put the right value in its place. And make sure you move left/right for X and up/down for Y.
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joeasanchez1223 4 years ago Posted 4 years ago. Direct link to joeasanchez1223's post “how do you know if the in...” more how do you know if the intercept you found is right? And why do we have to put the answer as an irrational? and how do we know when to put it as an irrational?
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David Severin 4 years ago Posted 4 years ago. Direct link to David Severin's post “You can always plug in th...” more You can always plug in the point to see if it works in the equation. So with -5x+4y=20, you have an x intercept of (-4,0) and -5(-4)+4(0)=20 gives 20=20, so it is correct. With (0,5), -5(0)+4(5)=20 also gives 20=20, so they work.
I am not sure what you mean by having the answer as an irrational number in linear equations. Almost all linear equations will have rational solutions especially what you deal with in Algebra I. It is more likely to get irrational solutions when you are talking about quadratic equations.
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CDOG55 2 years ago Posted 2 years ago. Direct link to CDOG55's post “This is very hard for me ...” more This is very hard for me my question is y= 10x - 32
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Brady Houser 2 years ago Posted 2 years ago. Direct link to Brady Houser's post “If you are trying to find...” more If you are trying to find the intercepts, from looking at the equation, you can see that the y-intercept is -32. To find the x-intercept, set y=0. That means you change y into a 0. So instead of y=10x-32, change it to 0=10x-32. Then add 32 to both sides of the equation, and then divide both sides by 10. Your x-intercept will be 3.2
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єɭɭค 2 years ago Posted 2 years ago. Direct link to єɭɭค's post “How do I solve an equatio...” more How do I solve an equation with a variable only on one side of the equation? Here's the question I'm on right now:
Determine the intercepts of the line
(Don't round your answer)
y-3=(x-2)
I get to solve for the intercepts you replace the x with a zero to get the y-intercept, but I ended up with this equation:
y-3=5(0-2)
Am I doing it right? And if I am, how do I solve the equation from here?
Please, somebody, anybody, I am confused and need help
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Kim Seidel 2 years ago Posted 2 years ago. Direct link to Kim Seidel's post “The work you have so far ...” more The work you have so far is ok to find the y-intercept.
1) Your next step is to simplify the right side. Follow order of operations rules (PEMDAS) to get to a single number on that side.
2) Then add 3 to each side to isolate "y".
Hope this helps.
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Video transcript
We have the equation negative 5x plus 4y is equal to 20, and we're told to find the intercepts of this equation. So we have to find the intercepts and then use the intercepts to graph this line on the coordinate plane. So then graph the line. So whenever someone talks about intercepts, they're talking about where you're intersecting the x and the y-axes. So let me label my axes here, so this is the x-axis and that is the y-axis there. And when I intersect the x-axis, what's going on? What is my y value when I'm at the x-axis? Well, my y value is 0, I'm not above or below the x-axis. Let me write this down. The x-intercept is when y is equal to 0, right? And then by that same argument, what's the y-intercept? Well, if I'm somewhere along the y-axis, what's my x value? Well, I'm not to the right or the left, so my x value has to be 0, so the y-intercept occurs when x is equal to 0. So to figure out the intercepts, let's set y equal to 0 in this equation and solve for x, and then let's set x is equal to 0 and then solve for y. So when y is equal to 0, what does this equation become? I'll do it in orange. You get negative 5x plus 4y. Well we're saying y is 0, so 4 times 0 is equal to 20. 4 times 0 is just 0, so we can just not write that. So let me just rewrite it. So we have negative 5x is equal to 20. We can divide both sides of this equation by negative 5. The negative 5 cancel out, that was the whole point behind dividing by negative 5, and we get x is equal to 20 divided by negative 5 is negative 4. So when y is equal to 0, we saw that right there, x is equal to negative 4. Or if we wanted to plot that point, we always put the x coordinate first, so that would be the point negative 4 comma 0. So let me graph that. So if we go 1, 2, 3, 4. That's a negative 4. And then the y value is just 0, so that point is right over there. That is the x-intercept, y is 0, x is negative 4. Notice we're intersecting the x-axis. Now let's do the exact same thing for the y-intercept. Let's set x equal to 0, so if we set x is equal to 0, we have negative 5 times 0 plus 4y is equal to 20. Well, anything times 0 is 0, so we can just put that out of the way. And remember, this was setting x is equal to 0, we're doing the y-intercept now. So this just simplifies to 4y is equal to 20. We can divide both sides of this equation by 4 to get rid of this 4 right there, and you get y is equal to 20 over 4, which is 5. So when x is equal to 0, y is equal to 5. So the point 0, 5 is on the graph for this line. So 0, 5. x is 0 and y is 1, 2, 3, 4, 5, right over there. And notice, when x is 0, we're right on the y-axis, this is our y-intercept right over there. And if we graph the line, all you need is two points to graph any line, so we just have to connect the dots and that is our line. So let me connect the dots, trying my best to draw as straight of a line is I can-- well, I can do a better job than that-- to draw as straight of a line as I can. And that's the graph of this equation using the x-intercept and the y-intercept.
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13745 | https://ultimatepopculture.fandom.com/wiki/Ox | Ox | Ultimate Pop Culture Wiki | Fandom
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Contents
1 Domestication
2 Training
3 Shoeing
4 Uses and comparison to horses
5 See also
6 References
7 External links
in:Articles with incorrect citation syntax, CS1 Italian-language sources (it), CS1 Sardinian-language sources (sc),
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Cattle
Pack animals
Working animals
Ox
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This article is about cattle used for draft. For other uses of ox or oxen , see Ox (disambiguation). For other uses of bullock , see Bullock.
Not to be confused with aurochs or musk ox.
Zebu oxen in Mumbai, India
File:Ploughing with Oxen.jpg
Ploughing with Oxen by George H. Harvey, Nova Scotia, Canada, 1881
File:Traditional Farming Methods and Equipments.jpg
Oxen used in farms for plowing
An ox (pluraloxen), also known as a bullock in Australia and India, is a bovine trained as a draft animal. Oxen are commonly castrated adult male cattle; castration makes the animals easier to control. Cows (adult females) or bulls (intact males) may also be used in some areas.
Oxen are used for plowing, for transport (pulling carts, hauling wagons and even riding), for threshing grain by trampling, and for powering machines that grind grain or supply irrigation among other purposes. Oxen may be also used to skid logs in forests, particularly in low-impact, select-cut logging.
Oxen are usually yoked in pairs. Light work such as carting household items on good roads might require just one pair, while for heavier work, further pairs would be added as necessary. A team used for a heavy load over difficult ground might exceed nine or ten pairs.
[x]
Contents
1 Domestication
2 Training
3 Shoeing
4 Uses and comparison to horses
5 See also
6 References
7 External links
Domestication[]
Oxen are thought to have first been harnessed and put to work around 4000 BC.
Training[]
File:Bullock wagon Promontory Road.jpg
A team of ten pair of oxen in Australia
Working oxen are taught to respond to the signals of the teamster, bullocky or ox-driver. These signals are given by verbal command and body language, reinforced by a goad, whip or a long pole (which also serves as a measure of length: see rod). In pre-industrial times, most teamsters were known for their loud voices and forthright language.
Verbal commands for draft animals vary widely throughout the world. In North America, the most common commands are:
Back: back up
Gee: turn to the right
Get up (also giddyup or giddyap, contractions for "get thee up" or "get ye up"): go
Haw: turn to the left
Whoa: stop
In the New England tradition, young castrated cattle selected for draft are known as working steers and are painstakingly trained from a young age. Their teamster makes or buys as many as a dozen yokes of different sizes for each animal as it grows. The steers are normally considered fully trained at the age of four and only then become known as oxen.
A tradition in south eastern England was to use oxen (often Sussex cattle) as dual-purpose animals: for draft and beef. A plowing team of eight oxen normally consisted of four pairs aged a year apart. Each year, a pair of steers of about three years of age would be bought for the team and trained with the older animals. The pair would be kept for about four years, then sold at about seven years old to be fattened for beef – thus covering much of the cost of buying that year's new pair. Use of oxen for plowing survived in some areas of England (such as the South Downs) until the early twentieth century. Pairs of oxen were always hitched the same way round, and they were often given paired names. In southern England it was traditional to call the near-side (left) ox of a pair by a single-syllable name and the off-side (right) one by a longer one (for example: Lark and Linnet, Turk and Tiger).
Ox trainers favor larger animals for their ability to do more work. Oxen are therefore usually of larger breeds, and are usually males because they are generally larger. Females can also be trained as oxen, but as well as being smaller, are often more valued for producing calves and milk. Bulls are also used in many parts of the world, especially Asia and Africa.
Shoeing[]
Working oxen usually require shoes, although in England not all working oxen were shod. Since their hooves are cloven, two shoes are required for each hoof, as opposed to a single horseshoe. Ox shoes are usually of approximately half-moon or banana shape, either with or without caulkins, and are fitted in symmetrical pairs to the hooves. Unlike horses, oxen are not easily able to balance on three legs while a farrier shoes the fourth. In England, shoeing was accomplished by throwing the ox to the ground and lashing all four feet to a heavy wooden tripod until the shoeing was complete. A similar technique was used in Serbia and, in a simpler form, in India, where it is still practiced. In Italy, where oxen may be very large, shoeing is accomplished using a massive framework of beams in which the animal can be partly or completely lifted from the ground by slings passed under the body; the feet are then lashed to lateral beams or held with a rope while the shoes are fitted.
Such devices were made of wood in the past, but may today be of metal. Similar devices are found in France, Austria, Germany, Spain, Canada and the United States, where they may be called ox slings, ox presses or shoeing stalls. The system was sometimes adopted in England also, where the device was called a crush or trevis; one example is recorded in the Vale of Pewsey. The shoeing of an ox partly lifted in a sling is the subject of John Singer Sargent's painting Shoeing the Ox, while A Smith Shoeing an Ox by Karel Dujardin shows an ox being shod standing, tied to a post by the horns and balanced by supporting the raised hoof.
Ox shoe.jpg
A single left-hand ox shoe of the type used for large Chianina oxen in Tuscany
Karel Dujardin - A Smith Shoeing an Ox.jpg
Karel Dujardin, 1622–1678: A Smith Shoeing an Ox
2008.04.18.VorrichtungZumBeschlagenVonOchsen.DorfmuseumMoenchhof.33.JPG
Ox shoeing sling in the Dorfmuseum of Mönchhof, Austria; a pair of ox shoes is attached to the near left column.
Uses and comparison to horses[]
File:Sixten.jpg
Riding an ox in Hova, Sweden
Oxen can pull heavier loads, and pull for a longer period of time than horses depending on weather conditions. On the other hand, they are also slower than horses, which has both advantages and disadvantages; their pulling style is steadier, but they cannot cover as much ground in a given period of time. For agricultural purposes, oxen are more suitable for heavy tasks such as breaking sod or plowing in wet, heavy, or clay-filled soil. When hauling freight, oxen can move very heavy loads in a slow and steady fashion. They are at a disadvantage compared to horses when it is necessary to pull a plow or load of freight relatively quickly.
For millennia, oxen also could pull heavier loads because of the use of the yoke, which was designed to work best with the neck and shoulder anatomy of cattle. Until the invention of the horse collar, which allowed the horse to engage the pushing power of its hindquarters in moving a load, horses could not pull with their full strength because the yoke was incompatible with their anatomy (yokes press on their chest, inhibiting their breathing).
Well-trained oxen are also considered less excitable than horses.
See also[]
Aurochs
Bullock cart (ox-cart)
Bullocky (ox-driver, teamster)
Ox (zodiac)
Ox in Chinese mythology
Ox-wagon (bullock wagon)
Oxtail
Ridge and furrow
References[]
↑"HISTORY OF THE DOMESTICATION OF ANIMALS". Historyworld.net. Archived from the original on November 24, 2012. Retrieved September 17, 2012.{{cite web}}:
↑Conroy, Drew (2007). Oxen, A Teamster's Guide. North Adams, Massachusetts, USA: Storey Publishing. ISBN978-1-58017-693-4.
↑Copper, Bob, A Song for Every Season: A Hundred Years of a Sussex Farming Family (pp 95–100), Heinemann 1971
↑John C Barret (1991), "The Economic Role of Cattle in Communal Farming Systems in Zimbabwe", to be published in Zimbabwe Veterinary Journal, p 10.Archived 2012-09-18 at the Wayback Machine
↑Draught Animal Power, an Overview, Agricultural Engineering Branch, Agricultural Support Systems Division, Food and Agriculture Organisation of the United NationsArchived 2010-07-01 at the Wayback Machine
↑ Jump up to: 6.06.16.2Williams, Michael (17 September 2004). "The Living Tractor". Farmers Weekly. Archived from the original on 3 March 2011. Retrieved 22 May 2011.
↑ Jump up to: 7.07.1Watts, Martin (1999). Working oxen. Princes Risborough: Shire. ISBN0-7478-0415-X. Archived on 2014-06-12. Error: If you specify |archivedate=, you must also specify |archiveurl=.
↑ Jump up to: 8.08.1Baker, Andrew (1989). "Well Trained to the Yoke: Working Oxen on the Village's Historical Farms". Old Sturbridge Village. Archived from the original on 26 September 2011. Retrieved 22 May 2011.{{cite web}}:
↑Schomberg, A. (7 November 1885). "Shoeing oxen and horses at a Servian smithy". The Illustrated London News. Archived from the original on 4 October 2011. Retrieved 22 May 2011.
↑"Blacksmith shoeing a Bullock, Calcutta, India"(stereoscope card (half only)). Stereo-Travel Co. 1908. Retrieved 22 May 2011.{{cite web}}:
↑Aliaaaaa (2006). "Restraining and Shoeing". Bangalore, Karnataka, India. Archived from the original on 20 December 2013. Retrieved 22 May 2011.{{cite web}}:
↑Tacchini, Alvaro. "La ferratura dei buoi" . Archived from the original on 11 December 2013. Retrieved 22 May 2011. The shoeing of the oxen{{cite web}}:
↑"Tradizioni - Serramanna" . Archived from the original on 7 October 2011. Retrieved 22 May 2011. Serramanna: traditions{{cite web}}:
↑"Did You Know?". Wet Dry Routes Chapter Newlsletter4 (4). 1997. Archived from the original on 22 July 2011. Retrieved 22 May 2011.
↑John Singer Sargent. "Shoeing the Ox". Archived from the original on 11 July 2016. Retrieved 18 August 2016.{{cite web}}:
↑Taylor, Tess (May 3, 2011). "On Small Farms, Hoof Power Returns". The New York Times. Archived from the original on 22 May 2013.
↑Conroy, Drew. "Dr"(PDF). Ox Yokes: Culture, Comfort and Animal Welfare. World Association for Transport Animal Welfare and Studies (TAWS). Archived(PDF) from the original on 22 March 2012. Retrieved 21 February 2012.{{cite web}}:
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Combined arithmetic and geometric progression problem
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Asked 10 years, 5 months ago
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Tom fertilizes his plant once a week. Every week the amount of total fertilizers reduces by 25%. Let's say that every Saturday Tom adds 10 grams of fertilizers to his plant.
Write an equation to calculate the amount of fertilizers after every fertilization.
So first I realized that
a n=a n−1⋅q+d a n=a n−1⋅q+d
which in this case is
a n=0.75 a n−1+10 a n=0.75 a n−1+10
a.k.a
(((10⋅0.75+10)⋅0.75+10)⋅0.75+10)⋯(((10⋅0.75+10)⋅0.75+10)⋅0.75+10)⋯
I have also noted that the amount of reduction of fertilizers is going to be 1 less than the amount of addition of fertilizers.
But how to write this as an equation? I haven't progressed for a while now with this problem and I hope you can help me with an explanation behind your thoughts!
geometric-progressions
arithmetic-progressions
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edited Nov 25, 2015 at 13:07
Tacet
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asked Apr 18, 2015 at 13:54
BenskeyBenskey
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Let a m=b m+C a m=b m+C
⟹b n+C=0.75[b n−1+C]+10⟹b n+C=0.75[b n−1+C]+10
⟺b n+C=0.75 b n−1+0.75 C+10⟺b n+C=0.75 b n−1+0.75 C+10
Set C=0.75 C+10⟺C=?C=0.75 C+10⟺C=? to get ⟺b n=0.75 b n−1=⋯(0.75)r b n−r⟺b n=0.75 b n−1=⋯(0.75)r b n−r where r r is any integer.
You must have an initial condition, I presume?
In effect we find, a n−40=0.75[a n−1−40]a n−40=0.75[a n−1−40]
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answered Apr 18, 2015 at 14:03
lab bhattacharjeelab bhattacharjee
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{a 0=k a n=a n−1 q+d{a 0=k a n=a n−1 q+d
Is solved by
{if q≠1,a n=q n k+q n−1 q−1 d if q=1,a n=k+n d{if q≠1,a n=q n k+q n−1 q−1 d if q=1,a n=k+n d
In your case, q=0.75,d=10 q=0.75,d=10 and k k being the initial quantity of fertilizer (which you didn't state).
Note that, after a long period of time (a.k.a. for big values of n n), since |q|<1|q|<1, a n∼d 1−q=10 0.25=40 a n∼d 1−q=10 0.25=40, idependently of the initial dose of fertilizer.
Added
For q≠1 q≠1, let's show that q n k+q n−1 q−1 d q n k+q n−1 q−1 d is the solution of
{a 0=k a n=a n−1 q+d{a 0=k a n=a n−1 q+d
(if we use the convention 0 0:=1 0 0:=1)
Indeed, let m m the smallest integer such that a m≠q m k+q m−1 q−1 d a m≠q m k+q m−1 q−1 d.
Then, m>0 m>0, since a 0=k a 0=k and q 0 k+q 0−1 q−1=k+0 q−1 q 0 k+q 0−1 q−1=k+0 q−1.
But m>0→m=t+1 m>0→m=t+1, with t≥0 t≥0.
By minimality of m m, a t=q t k+q t−1 q−1 d a t=q t k+q t−1 q−1 d. By definition,
a m=a t q+d=q t+1 k+q t+1−q q−1 d+d=q m k+q m−1 q−1 d a m=a t q+d=q t+1 k+q t+1−q q−1 d+d=q m k+q m−1 q−1 d
Which yields a contradiction.
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edited Apr 18, 2015 at 15:43
answered Apr 18, 2015 at 14:18
user228113 user228113
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The initial state isn't stated in the problem either, but I assume it's 0, and after the first fertiliation it is 10.Benskey –Benskey 2015-04-18 15:26:24 +00:00 Commented Apr 18, 2015 at 15:26
Could you explain your thoughts behind the line where q != 1 ?Benskey –Benskey 2015-04-18 15:28:56 +00:00 Commented Apr 18, 2015 at 15:28
Then k=0 k=0, hence a n=40⋅(1−0.75 n)a n=40⋅(1−0.75 n). Ok, I'll add a proof (the explanation would be that I know the formula for ∑n 0 α s∑0 n α s)user228113 –user228113 2015-04-18 15:29:29 +00:00 Commented Apr 18, 2015 at 15:29
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13747 | https://www.youtube.com/watch?v=UwiUVKdttlY | Find the limit as x approaches a of (x^3 - ax^2 + a^2x -a^3)/(x-a)
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Transcript:
hi everyone we're going to find the limit as x approaches a of x cubed minus ax squared plus a squared times x minus a cubed and then we're dividing by x minus a so clearly if you substitute in a in for x we're going to get undefined in the denominator so let's regroup this so we can get x minus a in the numerator and the way to regroup this is let's see we'll do the limit as x approaches a and let's group these two together so we can use this formula here so that's going to be x cubed minus a cubed and then we still have these two so let's factor out a ax and that would give me x minus a all right and i'll divide it by x minus a so the next step is let's expand this we have the oh it's not to infinity it's the a right all right so now we have the limit as x approaches a and expanding this we're going to get x minus a i'm using this formula here times x squared plus a x plus a squared and we still have minus a x divided by x or times x minus a all it by i'm going to split this up i'm going to put the x minus say under here and x minus a under here all right so what will happen is these two will cancel these two will cancel so let's see what we have limit as x minus as x approaches a of x squared plus ax plus a squared minus ax so that's just going to give me when i cancel those out x squared plus ax so once you have that let's go ahead and subset let's see if we can do it on the same page here so i'll do it right up here so now that we have this we can just plug the a into this and that's going to give me um let's see x squared plus this is not this is supposed to be uh not an x here okay i knew i did something wrong all right so when you substitute in your a you're going to get a squared plus a squared well that equals 2 a squared so 2a squared is your limit and a is some constant and that's it thank you have a nice day bye |
13748 | https://www.51jiaoxi.com/doc-17198429.html | 资料篮
人教课标B版高中数学必修3 2-2-1《用样本的频率分布估计总体分布》第二课时教学设计案
所属成套资源:人教课标B版高中数学必修第三册 课件+教案+试题
高中数学用样本的频率分布估计总体第二课时教学设计
这是一份高中数学用样本的频率分布估计总体第二课时教学设计,共2页。教案主要包含了复习准备,讲授新课,课堂小结等内容,欢迎下载使用。
1.通过实例体会分布的意义和作用,在表示样本数据的过程中,学会列频率分布表、画频率分布直方图、频率折线图、茎叶图,体会它们各自的特点.
2.在解决统计问题的过程中,进一步体会用样本估计总体的思想,会用样本的频率分布估计总体分布,
教学重点
学会列频率分布表、画频率分布直方图、频率折线图、茎叶图.
教学难点
体会用样本估计总体的思想,会用样本的频率分布估计总体分布
教学过程:
一、复习准备:
1.讨论:绘制频率分布直方图有哪几个步骤呢?
2.练习:给出一个频率分布直方图,进行一些分析.
(如何表示频率?面积和?集中范围?变化趋势?)
二、讲授新课:
教学频率分布折线图及茎叶图:
1.定义频率分布折线图
画好频率分布图后,我们把频率分布直方图中各小长方形上端连接起来,得到的图形.
2.定义总体密度曲线
在样本频率分布直方图中,相应的频率折线图会越来越接近于一条光滑曲线,统计中称这条光滑曲线为总体密度曲线. 它能够精确地反映了总体在各个范围内取值的百分比,它能给我们提供更加精细的信息.
注:频率折线图是随着样本而变化的,因此并不能由频率折线图得到准确的总体密度曲线. 当样本容量不断增加,分组的组距不断缩小,频率分布折线图会越来越接近一条光滑的曲线即总体密度曲线,它由(a,b)的阴影部分的面积,直观反映总体在范围(a,b)内取值的百分比.
3.讨论:
对于任何一个总体,它的密度曲线是不是一定存在?它的密度曲线是否可以被非常准确地画出来?
实际上,尽管有些总体密度曲线是客观存在的,但一般很难想函数图象那样准确地画出来,我们只能用样本的频率分布对它进行估计,一般来说,样本容量越大,这种估计就越精确.
4.提问:
目前有哪些方式可以发现样本的规律?
(分布表、直方图、折线图都能帮助发现样本数据的规律)
5.定义茎叶图:
当数据是两位有效数字时,用中间的数字表示十位数,即第一个有效数字,两边的数字表示个位数,即第二个有效数字,它的中间部分像植物的茎,两边部分像植物茎上长出来的叶子,因此通常把这样的图叫做茎叶图.
注:茎叶是一种形象的说法,表明两部分数据间的关系,茎是指数据中用来分组的依据数,叶是指被分到这组的数.
6.出示例:试将下列两组数据制作出茎叶图.
甲得分:13 ,51,23,8,26,38,16,33,14,25,39,
乙得分:49,24,12,31,60,31,44,36,15,37,25,36,39,
(▲ 师生共同按制作茎叶图的方法进行操作)
7.讨论:用茎叶图处理样本数据有何好处,什么时候用茎叶图会比较方使?
(茎叶图不仅能够保留原始数据,数据可以随时记录,随时添加,方便记录,而且能够展示数据的分布情况,但其仅适用于样本数据较少时,否则枝叶会太长.茎叶图中数据的茎和叶的划分,可根据数据的特点灵活地决定.)
三、课堂小结:
不易知一个总体的分布情况时,往往从总体中抽取一个样本,用样本的频率分布去估计总体的频率分布,样本容量越大,估计就越精确.目前有:频率分布表、直方图、茎叶图.
相关教案
这是一份高中数学人教版新课标B必修3用样本的频率分布估计总体第一课时教案,共3页。教案主要包含了复习准备,讲授新课等内容,欢迎下载使用。
这是一份高中数学人教版新课标A必修32.2.1用样本的频率分布估计总体教案及反思,共5页。教案主要包含了三维目标,重点与难点,教学设想等内容,欢迎下载使用。
这是一份2020-2021学年2.2.1用样本的频率分布估计总体教案及反思,共8页。教案主要包含了创设情境,探究新知,例题精析,课堂精练,课堂小结,评价设计等内容,欢迎下载使用。
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高中数学人教版新课标A必修32.2.1用样本的频率分布估计总体教案
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人教版新课标B必修32.2.1用样本的频率分布估计总体教案设计
人教版新课标A必修32.2.1用样本的频率分布估计总体教案设计
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Annals of Oncology
Volume 14, Issue 12, December 2003, Pages 1705-1713
Reviews The role of the E-cadherin gene (CDH1) in diffuse gastric cancer susceptibility: from the laboratory to clinical practice
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Abstract
Loss of function of the E-cadherin gene (CDH1) has been linked with diffuse gastric cancer susceptibility, and germline inactivating mutations in CDH1 characterise the hereditary diffuse gastric cancer (HDGC) syndrome. Hypermethylation in the CDH1 promoter region is a frequent phenomenon in poorly differentiated, diffuse gastric carcinomas and it was identified as the main mechanism for the inactivation of the remaining wild-type allele in HDGC cases. Specific criteria are used to identify patients with suspected HDGC and who should be investigated for CDH1 germline mutations. Accurate screening is mandatory for unaffected carriers of CDH1 mutations and selected high-risk individuals could be considered for prophylactic gastrectomy. Also, germline CDH1 mutations may predispose to lobular breast carcinoma and prostate cancer.Germline CDH1 mutations are not always detectable in patients who meet the HDGC criteria and the aetiological role of this gene is still under investigation. Families without recognised inactivating CDH1 mutations may have undisclosed CDH1 mutations or mutations in its regulatory sequences or germline mutations in unidentified genes that also contribute to the disease. In recent years, several germline missense CDH1 mutations have been identified, some of which showed a marked negative influence on E-cadherin function in experimental models. CDH1 promoter hypermethylation seems a key event in the carcinogenetic process of poorly differentiated, diffuse gastric cancer and it deserves further investigation as a new target for anticancer therapies with demethylating agents.
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E-cadherin
gastric cancer
germline
hereditary diffuse gastric cancer
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Copyright © 2003 European Society for Medical Oncology. Published by Elsevier Ltd. All rights reserved. |
13751 | https://phys.libretexts.org/Bookshelves/College_Physics/College_Physics_1e_(OpenStax)/31%3A_Radioactivity_and_Nuclear_Physics/31.05%3A_Half-Life_and_Activity | Skip to main content
31.5: Half-Life and Activity
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31.4: Nuclear Decay and Conservation Laws
31.6: Binding Energy
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OpenStax
OpenStax
( \newcommand{\kernel}{\mathrm{null}\,})
Learning Objectives
By the end of this section, you will be able to:
Define half-life.
Define dating.
Calculate age of old objects by radioactive dating.
Unstable nuclei decay. However, some nuclides decay faster than others. For example, radium and polonium, discovered by the Curies, decay faster than uranium. This means they have shorter lifetimes, producing a greater rate of decay. In this section we explore half-life and activity, the quantitative terms for lifetime and rate of decay.
Half-Life
Why use a term like half-life rather than lifetime? The answer can be found by examining Figure , which shows how the number of radioactive nuclei in a sample decreases with time. The time in which half of the original number of nuclei decay is defined as the half-life . Half of the remaining nuclei decay in the next half-life. Further, half of that amount decays in the following half-life. Therefore, the number of radioactive nuclei decreases from to in one half-life, then to in the next, and to in the next, and so on. If is a large number, then many half-lives (not just two) pass before all of the nuclei decay. Nuclear decay is an example of a purely statistical process. A more precise definition of half-life is that each nucleus has a 50% chance of living for a time equal to one half-life . Thus, if is reasonably large, half of the original nuclei decay in a time of one half-life. If an individual nucleus makes it through that time, it still has a 50% chance of surviving through another half-life. Even if it happens to make it through hundreds of half-lives, it still has a 50% chance of surviving through one more. The probability of decay is the same no matter when you start counting. This is like random coin flipping. The chance of heads is 50%, no matter what has happened before.
There is a tremendous range in the half-lives of various nuclides, from as short as s for the most unstable, to more than y for the least unstable, or about 46 orders of magnitude. Nuclides with the shortest half-lives are those for which the nuclear forces are least attractive, an indication of the extent to which the nuclear force can depend on the particular combination of neutrons and protons. The concept of half-life is applicable to other subatomic particles, as will be discussed in Particle Physics. It is also applicable to the decay of excited states in atoms and nuclei. The following equation gives the quantitative relationship between the original number of nuclei present at time zero (N_0) and the number () at a later time .
where is the base of the natural logarithm, and is the decay constant for the nuclide. The shorter the half-life, the larger is the value of and the faster the exponential decreases with time. The relationship between the decay constant and the half-life is
To see how the number of nuclei declines to half its original value in one half-life, let in the exponential in the equation . This gives
For integral numbers of half-lives, you can just divide the original number by 2 over and over again, rather than using the exponential relationship. For example, if ten half-lives have passed, we divide by 2 ten times. This reduces it to . For an arbitrary time, not just a multiple of the half-life, the exponential relationship must be used.
Radioactive dating is a clever use of naturally occurring radioactivity. Its most famous application is carbon-14 dating. Carbon-14 has a half-life of 5730 years and is produced in a nuclear reaction induced when solar neutrinos strike in the atmosphere. Radioactive carbon has the same chemistry as stable carbon, and so it mixes into the ecosphere, where it is consumed and becomes part of every living organism. Carbon-14 has an abundance of 1.3 parts per trillion of normal carbon. Thus, if you know the number of carbon nuclei in an object (perhaps determined by mass and Avogadro’s number), you multiply that number by to find the number of nuclei in the object. When an organism dies, carbon exchange with the environment ceases, and is not replenished as it decays. By comparing the abundance of in an artifact, such as mummy wrappings, with the normal abundance in living tissue, it is possible to determine the artifact’s age (or time since death). Carbon-14 dating can be used for biological tissues as old as 50 or 60 thousand years, but is most accurate for younger samples, since the abundance of nuclei in them is greater. Very old biological materials contain no at all. There are instances in which the date of an artifact can be determined by other means, such as historical knowledge or tree-ring counting. These cross-references have confirmed the validity of carbon-14 dating and permitted us to calibrate the technique as well. Carbon-14 dating revolutionized parts of archaeology and is of such importance that it earned the 1960 Nobel Prize in chemistry for its developer, the American chemist Willard Libby (1908–1980).
One of the most famous cases of carbon-14 dating involves the Shroud of Turin, a long piece of fabric purported to be the burial shroud of Jesus (Figure ). This relic was first displayed in Turin in 1354 and was denounced as a fraud at that time by a French bishop. Its remarkable negative imprint of an apparently crucified body resembles the then-accepted image of Jesus, and so the shroud was never disregarded completely and remained controversial over the centuries. Carbon-14 dating was not performed on the shroud until 1988, when the process had been refined to the point where only a small amount of material needed to be destroyed. Samples were tested at three independent laboratories, each being given four pieces of cloth, with only one unidentified piece from the shroud, to avoid prejudice. All three laboratories found samples of the shroud contain 92% of the found in living tissues, allowing the shroud to be dated (Example ).
Example : How Old Is the Shroud of Turin?
Calculate the age of the Shroud of Turin given that the amount of found in it is 92% of that in living tissue.
Strategy
Knowing that 92% of the remains means that . Therefore, the equation can be used to find . We also know that the half-life of is 5730 y, and so once is known, we can use the equation to find and then find as requested. Here, we postulate that the decrease in is solely due to nuclear decay.
Solution
Solving the equation for N/N_0) gives
Thus,
Taking the natural logarithm of both sides of the equation yields
so that
Rearranging to isolate gives
Now, the equation can be used to find for . Solving for and substituting the known half-life gives
We enter this value into the previous equation to find .
Discussion
This dates the material in the shroud to 1988–690 = a.d. 1300. Our calculation is only accurate to two digits, so that the year is rounded to 1300. The values obtained at the three independent laboratories gave a weighted average date of a.d. . The uncertainty is typical of carbon-14 dating and is due to the small amount of in living tissues, the amount of material available, and experimental uncertainties (reduced by having three independent measurements). It is meaningful that the date of the shroud is consistent with the first record of its existence and inconsistent with the period in which Jesus lived.
There are other forms of radioactive dating. Rocks, for example, can sometimes be dated based on the decay of . The decay series for ends with , so that the ratio of these nuclides in a rock is an indication of how long it has been since the rock solidified. The original composition of the rock, such as the absence of lead, must be known with some confidence. However, as with carbon-14 dating, the technique can be verified by a consistent body of knowledge. Since has a half-life of y, it is useful for dating only very old materials, showing, for example, that the oldest rocks on Earth solidified about years ago.
Activity, the Rate of Decay
What do we mean when we say a source is highly radioactive? Generally, this means the number of decays per unit time is very high. We define activity to be the rate of decay expressed in decays per unit time. In equation form, this is
where is the number of decays that occur in time . The SI unit for activity is one decay per second and is given the name becquerel (Bq) in honor of the discoverer of radioactivity. That is,
Activity is often expressed in other units, such as decays per minute or decays per year. One of the most common units for activity is the curie (Ci), defined to be the activity of 1 g of , in honor of Marie Curie’s work with radium. The definition of curie is
or decays per second. A curie is a large unit of activity, while a becquerel is a relatively small unit. . In countries like Australia and New Zealand that adhere more to SI units, most radioactive sources, such as those used in medical diagnostics or in physics laboratories, are labeled in Bq or megabecquerel (MBq).
Intuitively, you would expect the activity of a source to depend on two things: the amount of the radioactive substance present, and its half-life. The greater the number of radioactive nuclei present in the sample, the more will decay per unit of time. The shorter the half-life, the more decays per unit time, for a given number of nuclei. So activity should be proportional to the number of radioactive nuclei, , and inversely proportional to their half-life, . In fact, your intuition is correct. It can be shown that the activity of a source is
where is the number of radioactive nuclei present, having half-life . This relationship is useful in a variety of calculations, as the next two examples illustrate.
Example : How Great is the Activity in Living Tissue?
Calculate the activity due to in 1.00 kg of carbon found in a living organism. Express the activity in units of Bq and Ci.
Strategy
To find the activity using the equation , we must know and . The half-life of can be found in Appendix B, and was stated above as 5730 y. To find , we first find the number of nuclei in 1.00 kg of carbon using the concept of a mole. As indicated, we then multiply by (the abundance of in a carbon sample from a living organism) to get the number of nuclei in a living organism.
Solution
One mole of carbon has a mass of 12.0 g, since it is nearly pure . (A mole has a mass in grams equal in magnitude to found in the periodic table.) Thus the number of carbon nuclei in a kilogram is
So the number of nuclei in 1 kg of carbon is
Now the activity is found using the equation . Entering known values gives
or decays per year. To convert this to the unit Bq, we simply convert years to seconds. Thus,
or 250 decays per second. To express in curies, we use the definition of a curie,
Thus,
Discussion
Our own bodies contain kilograms of carbon, and it is intriguing to think there are hundreds of decays per second taking place in us. Carbon-14 and other naturally occurring radioactive substances in our bodies contribute to the background radiation we receive. The small number of decays per second found for a kilogram of carbon in this example gives you some idea of how difficult it is to detect in a small sample of material. If there are 250 decays per second in a kilogram, then there are 0.25 decays per second in a gram of carbon in living tissue. To observe this, you must be able to distinguish decays from other forms of radiation, in order to reduce background noise. This becomes more difficult with an old tissue sample, since it contains less , and for samples more than 50 thousand years old, it is impossible.
Human-made (or artificial) radioactivity has been produced for decades and has many uses. Some of these include medical therapy for cancer, medical imaging and diagnostics, and food preservation by irradiation. Many applications as well as the biological effects of radiation are explored in Medical Applications of Nuclear Physics, but it is clear that radiation is hazardous. A number of tragic examples of this exist, one of the most disastrous being the meltdown and fire at the Chernobyl reactor complex in the Ukraine (Figure ). Several radioactive isotopes were released in huge quantities, contaminating many thousands of square kilometers and directly affecting hundreds of thousands of people. The most significant releases were of , , , , , and . Estimates are that the total amount of radiation released was about 100 million curies.
Human and Medical Applications
Example : What Mass of Escaped Chernobyl?
It is estimated that the Chernobyl disaster released 6.0 MCi of into the environment. Calculate the mass of released.
Strategy
We can calculate the mass released using Avogadro’s number and the concept of a mole if we can first find the number of nuclei released. Since the activity is given, and the half-life of is found in Appendix B to be 30.2 y, we can use the equation to find .
Solution
Solving the equation for gives
Entering the given values yields
Converting curies to becquerels and years to seconds, we get
One mole of a nuclide has a mass of grams, so that one mole of has a mass of 137 g. A mole has nuclei. Thus the mass of released was
Discussion
While 70 kg of material may not be a very large mass compared to the amount of fuel in a power plant, it is extremely radioactive, since it only has a 30-year half-life. Six megacuries (6.0 MCi) is an extraordinary amount of activity but is only a fraction of what is produced in nuclear reactors. Similar amounts of the other isotopes were also released at Chernobyl. Although the chances of such a disaster may have seemed small, the consequences were extremely severe, requiring greater caution than was used. More will be said about safe reactor design in the next chapter, but it should be noted that Western reactors have a fundamentally safer design.
Activity decreases in time, going to half its original value in one half-life, then to one-fourth its original value in the next half-life, and so on. Since , the activity decreases as the number of radioactive nuclei decreases. The equation for as a function of time is found by combining the equations and , yielding
where is the activity at . This equation shows exponential decay of radioactive nuclei. For example, if a source originally has a 1.00-mCi activity, it declines to 0.500 mCi in one half-life, to 0.250 mCi in two half-lives, to 0.125 mCi in three half-lives, and so on. For times other than whole half-lives, the equation must be used to find .
PHET EXPLORATIONS: ALPHA DECAY
Watch alpha particles escape from a polonium nucleus, causing radioactive alpha decay. See how random decay times relate to the half life.
Summary
Half-life is the time in which there is a 50% chance that a nucleus will decay. The number of nuclei as a function of time is where is the number present at , and is the decay constant, related to the half-life by
One of the applications of radioactive decay is radioactive dating, in which the age of a material is determined by the amount of radioactive decay that occurs. The rate of decay is called the activity :
The SI unit for is the becquerel (Bq), defined by
is also expressed in terms of curies (Ci), where
The activity of a source is related to and by
Since has an exponential behavior as in the equation , the activity also has an exponential behavior, given by where is the activity at .
Glossary
becquerel
: SI unit for rate of decay of a radioactive material
half-life
: the time in which there is a 50% chance that a nucleus will decay
radioactive dating
: an application of radioactive decay in which the age of a material is determined by the amount of radioactivity of a particular type that occurs
decay constant
: quantity that is inversely proportional to the half-life and that is used in equation for number of nuclei as a function of time
carbon-14 dating
: a radioactive dating technique based on the radioactivity of carbon-14
activity
: the rate of decay for radioactive nuclides
rate of decay
: the number of radioactive events per unit time
curie
: the activity of 1g of , equal to
31.4: Nuclear Decay and Conservation Laws
31.6: Binding Energy |
13752 | https://www.ck12.org/flexi/cbse-math/higher-order-metric-conversions/convert-30-minutes-to-seconds/ | Convert 30 minutes to seconds. - Steps | CK-12 Foundation
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Convert 30 minutes to seconds.
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30 minutes is 1800 seconds.
Time conversion involves understanding the relationship between hours and minutes. Since 1 minute is equivalent to 60 seconds, converting minutes to seconds is a straightforward multiplication process.
Steps to convert minutes to seconds
Step 1: Identify the given time: The time provided is 30 minutes.
Step 2: Convert the minutes to seconds: Multiply the number of minutes by 60 to get the equivalent seconds:
30 × 60 = 1800 seconds.
Final Answer: 30 minutes is 1800 seconds.
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13753 | https://www.clevelandmetroschools.org/cms/lib/OH01915844/Centricity/Domain/274/Tables%20of%20Values.pdf | NAME CLASS DATE 36 SpringBoard® Course 3 Math Skills Workshop Unit 2 • Getting Ready Practice Tables of Values Tables of values are used to show the relationship between data items. You might use a table of values in science class. Scientists and researchers use tables of values to record their data and then analyze the data for a pattern. They can then use this pattern to make predictions. For example, a geologist who is studying earthquake activity might use a table of values to track the amount of underground movement versus distance from a fault line. She could then use the table to predict the distance the movement can be felt.
To complete a table of values that has missing numbers, look for a pattern in the data.
Complete the table of values. x y 23 17 21 14 1 3 8 5 7 2 Step 1: Look for a pattern in the left column.
Step 2: Use the pattern rule to find the missing values in the left column.
Step 3: Look for a pattern in the right column.
Step 4: Use the pattern rule to find the missing values in the right column.
Solution: x y 23 17 21 14 1 11 3 8 5 5 7 2 9 21 The numbers are increasing by 2.
3 1 2 5 5 7 1 2 5 9 The numbers are decreasing by 3.
14 2 3 5 11 2 2 3 5 21 EXAMPLE A 37 SpringBoard® Course 3 Math Skills Workshop Unit 2 • Getting Ready Practice Tables of Values (continued) You can create a table of values from a linear equation by substituting values for x into the equation and solving for y.
Complete the table of values for the equation. y 5 3x 2 4 x y 22 21 0 1 Substitute each x-value into the equation and solve for y.
y 5 3(22) 2 4 y 5 3(21) 2 4 y 5 3(0) 2 4 y 5 3(1) 2 4 y 5 26 2 4 y 5 23 2 4 y 5 0 2 4 y 5 3 2 4 y 5 210 y 5 27 y 5 24 y 5 21 Solution: y 5 3x 2 4 x y 22 210 21 27 0 24 1 21 EXAMPLE B 38 SpringBoard® Course 3 Math Skills Workshop Unit 2 • Getting Ready Practice Tables of Values (continued) To determine whether an equation correctly represents data in a table of values, substitute each pair of x- and y-values into the equation and simplify. If you get a true statement in every case, the equation represents the data.
Which equation represents the data in the table? A. y 5 3x 2 3 C. y 5 4x 2 1 B. y 5 x 1 7 D. y 5 2x 2 5 Step 1: Substitute x- and y-values into the first equation. If any statement is false, move on to the next equation.
Step 2: Substitute x- and y-values into the second equation. If any statement is false, move on to the next equation.
Step 3: Substitute x- and y-values into the third equation. If any statement is false, move on to the next equation.
Step 4: Substitute x- and y-values into the fourth equation. If any statement is false, move on to another equation.
Solution: Equation D, y 5 2x 2 5, represents the data in the table. A. y 5 3x 2 3 29 0 3(22) 2 3 1 0 3(3) 2 3 29 0 29 (true) 1 0 6 (false) B. y 5 x 1 7 29 0 22 1 7 29 0 5 (false) C. y 5 4x 2 1 29 0 4(22) 2 1 1 0 4(3) 2 1 29 0 29 (true) 1 0 11 (false) D. y 5 2x 2 5 29 0 2(22) 2 5 11 0 2(8) 2 5 29 0 29 (true) 11 0 11 (true) 1 0 2(3) 2 5 21 0 2(13) 2 5 1 0 1 (true) 21 0 21 (true) EXAMPLE C PRACTICE Complete each table of values. 1. x y 24 4 21 8 12 5 8 20 2. x y 22 22 3 6 8 14 22 18 3. x y 210 4 28 22 28 24 214 22 x 22 3 8 13 y 29 1 11 21 39 SpringBoard® Course 3 Math Skills Workshop Unit 2 • Getting Ready Practice Complete the table of values for each equation. 4. y 5 2x 1 3 x y 22 21 0 1 2 5. y 5 x 2 6 x y 22 21 0 1 2 6. y 5 4x 2 2 x y 22 21 0 1 2 7. Which equation represents the data in the table? Show your work. x 2 4 6 8 y 1 5 9 13 A. y 5 1 2x B. y 5 x 21 C. y 5 2x 2 3 D. y 5 3x 2 5 8. Dave and Bonnie each made a table of values for the equation y 5 1 2x 2 2. Dave said his table is correct. Bonnie said her table is correct. Who is correct? Explain your answer. Dave’s table x 24 22 0 2 y 24 23 22 21 Bonnie’s table x 4 8 12 16 y 0 2 4 6 Tables of Values (continued) |
13754 | https://math.stackexchange.com/questions/1141555/how-do-i-solve-an-inequality-with-2-inverse-trigonometrical-functions-involved | algebra precalculus - How do I solve an inequality with 2 inverse trigonometrical functions involved? - Mathematics Stack Exchange
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How do I solve an inequality with 2 inverse trigonometrical functions involved?
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I haven't worked with this in a long time! All I remember is that increasing vs decreasing functions have the power of modifying the symbol.
arcsin(2 x)>arccos(2 x)arcsin(2 x)>arccos(2 x)
If I take the sin of the arcsin I will have 2 x 2 x, I am very confused on how to approach this, can you give me some insight, please?
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edited Feb 10, 2015 at 3:31
Michael Hardy
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asked Feb 10, 2015 at 3:06
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Use arccos(y)=π 2−arcsin(y)arccos(y)=π 2−arcsin(y) to get
2 arcsin 2 x>π 2⟺arcsin 2 x>π 4 2 arcsin2 x>π 2⟺arcsin2 x>π 4
As arcsin arcsin is increasing, we need 2 x>sin π 4=1 2–√2 x>sinπ 4=1 2
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answered Feb 10, 2015 at 4:29
lab bhattacharjeelab bhattacharjee
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Take sin of both sides, then use the fact that
sin cos−1 x=1−x 2−−−−−√sincos−1x=1−x 2
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answered Feb 10, 2015 at 4:04
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0Trigonometric inequation sin−1 2 x>cos−1 2 x sin−12 x>cos−12 x
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13755 | https://www.youtube.com/watch?v=EZ5tU45Ti_g | Plant Growth: Auxins and Gibberellins | Plants | Biology | FuseSchool
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Plant Growth: Auxins and Gibberellins | Plants | Biology | FuseSchool
If a plant has enough water, minerals and energy, it will grow, right? Well, sort of… but there is more to it - like why do plants bend towards the light and not just grow straight? And why does the stem grow up but the roots grow down? It isn’t as if a plant has eyes to tell it where the sun is.
Plants are packed full of hormones, sending messages around to its different parts. Where humans have the creatively named ‘growth hormone’, plants have hormones called auxins.
Auxin is produced in the stem tips and roots, and controls the direction of growth in response to different stimuli including light and gravity. Having been made in the tips of the stems and roots, auxin is moved in solution by diffusion to older parts of the plant. In the stem, the auxin causes the cells to change in elasticity. More elastic cells absorb more water, and can grow longer.
Strangely, though, stems and roots respond differently to high concentrations of auxins. Whilst the stem cells grow more, the root cells actually grow less. So auxins make plants grow, but why do they bend towards the light? How do they know to do this when they don’t have eyes?
The bending happens because the light hits the one side more and breaks down the auxins in that side of the stem. So then growth slows down on the ‘light’ side. The faster growth on the ‘dark’ side causes the shoots and leaves to turn towards the light - which is ideal for the plant for photosynthesis.
Auxin is produced in the tips of growing shoots. If the tips are cut off, then no auxin can be produced and so no plant growth. If the tips are covered, whilst auxin is still produced, light cannot break it down and so phototropism cannot occur: the plant just grows straight up and does not bend towards the light.
Auxins have the opposite effect on root cells. In roots, auxins cause less growth. The shaded side of roots contain more auxins, and so they grow less. This enables the ‘light’ side of the roots to grow more and bend away from the light.
And if that wasn’t weird enough, we have opposites happening with auxins and gravity too. In a horizontal root, the bottom side contains more auxins and grows less, so the root bends downwards in the direction of gravity. So positive geotropism. But of course, the stem responds differently.
In a horizontal stem, again the bottom side contains more auxins because it is less directly hit by sunlight. But because auxins cause growth in stems, the bottom side grows more causing the stem to bend upwards, against the direction of gravity, so negative geotropism.
But I am giving auxins too much credit; they don’t work alone. They have a partner in crime; cytokinins. You don’t need to know anything about these hormones other than the fact that they work alongside auxins. There is another plant hormone that you do need to be aware of… Gibberellins.
Once a seed germinates, the roots and shoots start to grow. But for this, the seed needs energy. Luckily, the seed releases a hormone called gibberellin which causes the starch in the seed to turn into sugars and provide the seed with energy to grow. As well as causing shoot growth, gibberellins can also stimulate flowering and fruits in some plants. And they also work with auxins to cause stem elongation.
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if a plant has enough water minerals and energy it will grow right well sort of but there is more to it like why do plants bend towards the light and not just grow straight and why is it that the stem grows up but the roots grow down it isn't as if the plant has eyes to tell it where the sun is you should already know about tropisms so phototropism and geotropism but if you have forgotten watch our video called plant hormones tropisms first but how does it actually do this having been made in the tips of the stems and roots auxin is moved in solution by diffusion to older parts of the plant in the stem the auxin causes the cells to change in elasticity more elastic cells absorb more water and can grow longer strangely though stems and roots respond differently to high concentrations of auxin whilst the stem cells grow more the root cells actually grow less so auxins make plants grow but why do they bend towards the light how do they know to do this when they don't have eyes the bending happens because the light hits on one side more and breaks down the auxins on that side of the stem so then growth slows down on the light side the faster growth on the dark side causes the shoots and leaves to turn towards the light which is ideal for the plant for photosynthesis auxin is produced in the tips of growing shoots if the tips are cut off then no auxin can be produced and so no plant growth if the tips are covered whilst auxin is still produced light cannot break it down and so phototropism cannot occur the plant just grows straight up and does not bend towards the light pretty straightforward so far now let's switch it up a bit auxins have the opposite effect on root cells in roots auxins cause less growth the shaded side of roots contain more auxin and so they grow less this enables the light side of the roots to grow more and bend away from the light and if that wasn't weird enough we have the opposite happening with auxins to gravity too in a horizontal route the bottom side contains more auxin and grows less so the root bends downwards in the direction of gravity so positive geotropism but of course the stem responds differently in a horizontal stem again the bottom side contains more auxins because it is less directly hit by sunlight but because auxins cause growth in stems the bottom side grows more causing the stem to bend upwards against the direction of gravity so negative geotropism but i'm giving auxins too much credit they don't work alone they have a partner in crime cytokinins you don't need to know anything about these hormones other than the fact that they work alongside auxins there is another plant hormone however that you do need to be aware of gibberellins once the seed germinates the roots and shoots start to grow but for this the seed needs energy luckily the seed releases a hormone called gibberellin which causes the starch in the seed to turn into sugars and provides the seed with energy to grow as well as causing shoot growth gibberillins can also stimulate flowering and the production of fruit in some plants and they also work with auxins to cause stem elongation so now we know all about how plant hormones work have a think about how they could be used we will discover some uses in our video on how we use plant hormones |
13756 | https://www.quora.com/What-is-an-intuitive-explanation-for-the-change-of-the-base-formula-for-logarithms | Something went wrong. Wait a moment and try again.
Change of Basis
Intuitive Explanations
Formula Transformation
Logarithmic Functions
Algebraic Concepts
Intuitive Solution
Changing Formulas
Intuitive Sense
5
What is an intuitive explanation for the change of the base formula for logarithms?
Dean Rubine
Former Faculty at Carnegie Mellon School Of Computer Science (1991–1994)
·
Author has 10.5K answers and 23.4M answer views
·
7y
logbx=logaxlogab
I always find this confusing and have a hard time writing it down from memory. But I recreate it like this:
Those are easy and self evident. Substituting,
Taking logarithms base of both sides,
That’s one form of the change of base. Alternatively,
I’ll admit that’s not all that intuitive. I can usually sort it out by remembering that the smaller the base the larger the magnitude of the logarithm of the sam
I always find this confusing and have a hard time writing it down from memory. But I recreate it like this:
Those are easy and self evident. Substituting,
Taking logarithms base of both sides,
That’s one form of the change of base. Alternatively,
I’ll admit that’s not all that intuitive. I can usually sort it out by remembering that the smaller the base the larger the magnitude of the logarithm of the same quantity.
So for example if , then
This helps us figure out what goes in the denominator, or Since the left side is bigger above, on the right we need a denominator smaller than one, so that must be :
That means
Philip Lloyd
Specialist Calculus Teacher, Motivator and Baroque Trumpet Soloist.
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Author has 6.8K answers and 52.5M answer views
·
7y
This indeed can be quite confusing.
I find so many people who don’t seem to grasp one very simple idea…
“Log” just means “Index”
I always feel sad that so many people actually NEED to consult a formula sheet in order to change the base of a logarithm!
The main reason for changing the base of a logarithm is that calculators only have logs to the base of 10 and e.
This indeed can be quite confusing.
I find so many people who don’t seem to grasp one very simple idea…
“Log” just means “Index”
I always feel sad that so many people actually NEED to consult a formula sheet in order to change the base of a logarithm!
The main reason for changing the base of a logarithm is that calculators only have logs to the base of 10 and e.
David Smith
BSc (Hons) in Mathematics & Computer Science, University of Bristol (Graduated 1986)
·
Author has 3.6K answers and 4M answer views
·
7y
Assume I can calculate logarithms to the base .
Thus, given a value , I can calculate by using .
Now: How to calculate ?
We need to find such that .
If we write for some then we have:
.
But we can take so we have .
Our value is defined such that .
So finally:
Related questions
What is the fundamental concept behind logarithms?
What is an intuitive explanation for exponential and logarithm curves?
What is the intuition behind the logarithm?
What is the mathematical meaning of a natural logarithm?
What is a logarithmic basis?
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The change of base formula for logarithms states that for any positive numbers , , and (where and ):
Intuitive Explanation
Understanding Logarithms:
A logarithm answers the question: "To what power must we raise the base to get ?" In mathematical terms, means .
Using a New Base:
Sometimes, it’s more convenient to express logarithms in a different base (like base 10 or base ). The change of base formula allows us to convert from one logarithmic base to another.
Breaking it Down:
Suppose you want to
The change of base formula for logarithms states that for any positive numbers , , and (where and ):
Intuitive Explanation
Understanding Logarithms:
A logarithm answers the question: "To what power must we raise the base to get ?" In mathematical terms, means .
Using a New Base:
Sometimes, it’s more convenient to express logarithms in a different base (like base 10 or base ). The change of base formula allows us to convert from one logarithmic base to another.
Breaking it Down:
Suppose you want to calculate . You can express in terms of base :
Let (so )
Let (so )
Now, you can rewrite in terms of base :
Since , we substitute to get or .
Equating Powers:
Since and , we can equate the exponents:
Dividing both sides by (assuming ), we obtain:
Thus, we can express as:
Summary
The change of base formula is a powerful tool that allows us to convert logarithms from one base to another by leveraging the relationship between the bases and the numbers involved. This flexibility is particularly useful in calculations or when using calculators that have specific bases (like base 10 or base ).
Aman Karunakaran
I read a lot about Math.
·
Author has 164 answers and 961.2K answer views
·
7y
From
Aman Karunakaran's answer to Why is
\log_a{b} = \frac{\log{b}}{\log{a}}
?
:
If [math]a^x = b[/math], then taking [math]\log_a[/math] of both sides we have [math]x = \log_a{b}[/math]
However, if we take the [math]\log[/math] of both sides we have [math]\log{a^x} = \log{b}[/math]
[math]\implies x\log{a} = \log{b}[/math]
[math]\implies x = \dfrac{\log{b}}{\log{a}}[/math]
And if we substitute [math]x[/math] from the beginning, we have:
[math]\log_a{b} = \dfrac{\log{b}}{\log{a}}[/math]
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There are a lot of AI tools out there right now—so how do you know which ones are actually worth your time? Which tools are built for students and school—not just for clicks or content generation? And more importantly, which ones help you sharpen what you already know instead of just doing the work for you?
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Ayodeji Oyenaike
Knows Japanese
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Author has 152 answers and 294.3K answer views
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7y
Let’s start with:
[math]\log_a(x) = y[/math]
Rearrange to an exponential equation:
[math]a^y = x[/math]
Take the [math]\log_b()[/math] of both sides:
[math]\log_b(a^y) = \log_b(x)[/math]
Solve for y:
[math]y \log_b(a) = \log_b(x)[/math]
[math]y = \dfrac{\log_b(a)}{\log_b(x)}[/math]
[math]\boxed{\therefore \log_a(x) = \dfrac{\log_b(a)}{\log_b(x)}}[/math]
Related questions
Why is a natural logarithm called so?
What do logarithms mean in layman's terms?
What is an easy and practical way to understand or explain logarithms and logarithmic scales?
What is a logarithm (with examples)?
What is the proper explanation of the concept of logarithms?
Drake Way
Mathematics Hobbyist
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·
6y
Originally Answered: What's the derivation of the logarithm change of base rule?
·
[math]c=\dfrac {\log_a x}{\log_a b}[/math]
[math]=(\log_a x)\dfrac 1 {\log_a b}[/math]
[math]=\log_a x^{\frac 1 {\log_a b}}[/math]
[math]a^c=x^{\frac 1 {\log_a b}}[/math]
math^{\log_a b}=x[/math]
[math]a^{(c\log_a b)}=x[/math]
math^c=x[/math]
[math]b^c=x[/math]
[math]c=\log_b x[/math]
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Related
What is the intuition behind the logarithm?
A lot of people say that a logarithm is an exponent. I technically disagree with that. Logarithms are functions, not numbers or values. But I get what people mean: the output of a logarithmic expression is an exponent. And I agree with that. So maybe I’m a bit pedantic. This isn’t the first time I’ve been told that.
But I don’t get how people can say that a “logarithm is an exponent” and then follow that statement up with “so [math]\log_2{9}[/math] is the power that I raise 2 to, to get 9”.
Okay, I’ve just conceded that a logarithm is an exponent (through gently gritted teeth) and now you want to call an exp
A lot of people say that a logarithm is an exponent. I technically disagree with that. Logarithms are functions, not numbers or values. But I get what people mean: the output of a logarithmic expression is an exponent. And I agree with that. So maybe I’m a bit pedantic. This isn’t the first time I’ve been told that.
But I don’t get how people can say that a “logarithm is an exponent” and then follow that statement up with “so [math]\log_2{9}[/math] is the power that I raise 2 to, to get 9”.
Okay, I’ve just conceded that a logarithm is an exponent (through gently gritted teeth) and now you want to call an exponent a power?
This is all too much for me.
Exponents are not powers.
Powers are not exponents.
You cannot replace the word power for the word exponent. Punto. This is our legacy, folks.
What are the powers of 10?
10, 100, 1000, 10000, … (among others) and can also be written as [math]10^1[/math], [math]10^2[/math], [math]10^3[/math], [math]10^4[/math], …
Did you notice that we didn’t say that the powers of 10 are 1, 2, 3, 4, …?
The powers of 10 are the output numbers (10, 100, 1000, 10000, …) of raising the base of 10 to the input EXPONENTS of 1, 2, 3, 4, … (or whatever exponent to which you want to raise the base).
Grammar, people.
The sons of Tom are José, Ahmed, Liam, [math]\theta[/math] and Vijay, so the third son of Tom is Liam. Not 3. Liam.
The multiples of 2 are 2, 4, 6, 8, 10, …, so the third multiple of 2 is 6. Not 3. Six.
The powers of 2 are 2, 4, 8, 16, 32, …, so the third power of 2 is 8. Not 3. Eight.
Okay. I feel a bit better. So now that we know that in [math]2^3=8[/math] that the base is 2, the exponent is 3 and the power is 8 (not 3), we can finally develop some intuition about logarithms.
The Logarithmic Question
The logarithmic expression [math]\log_b{p}[/math] asks “what is the exponent to which you raise the base of b to get the power of p?”. (The answer, of course is [math]\log_b{p}[/math]. Just like we might ask “what is the square root of 10?” and the answer is [math]\sqrt{10}[/math].)
For example, [math]\log_2{8}[/math] asks “what is the exponent to which you raise the base of 2 to get the power of 8?” The answer is 3 because [math]2^3=8[/math]. Sometimes you can simplify further and sometimes, you can’t:
[math]\sqrt{9}=3[/math] and [math]\log_2{8} = 3[/math]
[math]\sqrt{10} = \sqrt{10}[/math] and [math]\log_2{9} = \log_2{9}[/math] Aren’t these ones easier?
Exponential Form [math]\textrm{base}^\textrm{exponent}=\textrm{power}[/math] versus
Logarithmic Form [math]\log_{\textrm{base}}{\textrm{power}}=\textrm{exponent}[/math]
Both forms state the base, the exponent and the power.
There is a unique triad of base, exponent and power: [math]2^3[/math] is only 8 and [math]\log_2{8}[/math] is only 3.
In exponential form, the power is isolated. In logarithmic form, the exponent is isolated. These two equations convey the same three values but are mere rearrangements of each other.
By the way, if you want to know how we isolate the base, you need to use roots (the cube root in the case of [math]2^3=8[/math]). We’re not doing that here. We are only looking at the pair of inverse functions like 2^ and [math]\log_2[/math] (exponential and logarithmic functions, base 2) where the base of 2 is part of the function and NOT part of either the input or the output.
Exponent Laws and Log Laws Say the Same Thing … Just in a Different Way
Let [math]b^m=p\Longleftrightarrow \log_b{p}=m[/math] and [math]b^n=q\Longleftrightarrow \log_b{q}=n[/math].
[math]{b^m}\cdotp{b^n}=b^{m+n}[/math] and [math]\log_b{({p}\cdotp{q})}= \log_b{p}+ \log_b{q}[/math] both say (for the same base throughout) that adding exponents can be converted into multiplying powers. Notice however that when the exponents are added, the exponential form is condensed but the logarithmic form is expanded, and vice versa for when the powers are multiplied.
You can already guess where this is going. For the same base throughout, subtracting exponents can be converted into dividing powers: [math]{b^m}\div{b^n}=b^{m-n}[/math] and [math]\log_b{({p}\div{q})}= \log_b{p}- \log_b{q}[/math]. The same note applies about condensing and expanding.
The rest of the intuition follows from giving each element its own unique name. If we incorrectly say “two to the power 3 is 8” or “two to the third power is 8” (shudder; that’s even worse … do you mean an exponent of [math]\frac{1}{3}[/math]?!), then we will completely lose track of who is whom once we reflect through [math]y=x[/math] to get logarithmic functions, the inverse of exponential functions.
So to develop your intuition with logarithms, know your powers of 2, 3, 4, 5, 6, 10 (like you know your multiplication tables), know your exponent laws cold, understand the relationship between any function and its inverse, and
Do. Not. Ever.
call an exponent a power again.
Raelene Maths
Former AP Calculus, AP Physics, Head Teacher at American School of Doha (2004–2010)
·
Author has 63 answers and 180.7K answer views
·
Updated 6y
Related
Why can I not understand logarithms?
The main reason that people struggle to understand logarithms is due to the misuse of the word power.
Consider this sentence: The third son of Tom is Liam.
Who is the son? The son is Liam. Liam is the third son of Tom. Liam is the son, everything else in that sentence (the third, of Tom) just describes son.
Now consider this sentence, taken from [math]2^3=8[/math] and the “start at 2 and double” sequence {2, 4, 8, 16, 32, …} : The third power of 2 is 8.
Which number is the power? The power is 8. Eight is the third power of 2. 8 is the power, everything else is just a description of power.
Go ahead and search
“Po
The main reason that people struggle to understand logarithms is due to the misuse of the word power.
Consider this sentence: The third son of Tom is Liam.
Who is the son? The son is Liam. Liam is the third son of Tom. Liam is the son, everything else in that sentence (the third, of Tom) just describes son.
Now consider this sentence, taken from [math]2^3=8[/math] and the “start at 2 and double” sequence {2, 4, 8, 16, 32, …} : The third power of 2 is 8.
Which number is the power? The power is 8. Eight is the third power of 2. 8 is the power, everything else is just a description of power.
Go ahead and search
“Powers of 2” and you’ll get this list of numbers: {…, 2, 4, 8, 16, 32, …} NOT {…, 1, 2, 3, 4, 5, …};
or search “Powers of 10” and you’ll get {…, 10, 100, 1000, 10000, 100000, …} NOT {…, 1, 2, 3, 4, 5, …}.
If we call the parts of an exponential equation by the correct names, we will understand logarithms.
In the exponential form of this equation, [math]2^3=8[/math], 2 is the base, 3 is the exponent and 8 is the power. It is incorrect to use the word “power” as a synonym for “exponent”. An “index” is a special type of exponent; exponents are any real values while indices are natural (counting number) exponents.
[math]\textrm{base}^\textrm{exponent} = \textrm{power}[/math] is clear
[math]\textrm{base}^\textrm{power} = \textrm{power}[/math] (or result) is confusing and incorrect
We can rearrange equations to isolate different values, like
from addition to subtraction : [math]3 + 2 = 5 \Longleftrightarrow 5 - 2 = 3[/math]
from multiplication to division: [math]3 \times 2 = 6 \Longleftrightarrow 6 \div 2 = 3[/math]
from exponentiation to logarithmication (okay, I made this term up): [math]2^3=8 \Longleftrightarrow\log_{2}{8}=3[/math].
Same, same but different
Exponential form [math]2^3=8[/math] says: the POWER that results from the base 2 raised to the exponent 3 is 8.
Logarithmic form [math]\log_{2}{8}= 3[/math] says: the EXPONENT on the base of 2 to get the power of 8 is 3.
(Poor old base 2; he is just not the star of the exponential or logarithmic shows …)
Both forms say THE SAME THING.
The base is 2, the exponent is 3 and the power is 8.
These three values form a unique base-exponent-power triad, no matter how you rearrange them. Try it with these numbers: 100, 2 and 10 … (base = 10, exponent = 2 and power = 100: [math]10^2=100[/math]).
But they say it DIFFERENTLY.
The result of “raising a base to an exponent” is a power. In the exponential equation [math]2^3=8[/math], the power 8 is isolated.
The result of “taking the log for a given base of a power” is an exponent. In the logarithmic equation [math]\log_{2}{8}=3[/math], the exponent 3 is isolated.
How to Master Logarithms
Know basic number facts cold. Memorise some basic number facts (perfect squares, perfect cubes, perfect powers of 2, 3, 4, 5, 6, 10). (See tables below.)
Be more verbal. Start saying “two to the exponent three” for [math]2^3[/math] instead of “two to the power 3” or “two to the third power”. Repeat/emphasise the words “base”, “exponent” and “power” when reading an exponential or logarithmic expression or equation.
Translate to your “native” language. Rearrange logarithmic equations into exponential equations. When asked to evaluate the log expression [math]\log_{2}{32}[/math], first let the expression equal a variable so you have an equation to rearrange into exponential equation form:
Let [math]w=\log_{2}{32}[/math]. Now we have an exponent of w, a base of 2 and a power of 32.
Translate to exponential form: [math]2^w=32[/math].
Recall that [math]2^5=32[/math]. So, the exponent that you put on the base of 2 to get the power of 32 is 5.
Hence, w, the unique exponent for the base of 2 and power of 32 is 5: [math]\log_{2}{32}= 5[/math]. This sentence literally says that the exponent that you put on the base of 2 to get the power of 32 is 5. But in a much shorter format!
A Few More Examples
[math]\log_{2}{32}=5[/math] because [math]2^5=32[/math]
[math]\log_{7}{1}=0[/math] because [math]7^0=1[/math]
[math]\log_{3}{\frac{1}{9}}={-2}[/math] because [math]3^{-2}=\frac{1}{9}[/math]
[math]\log_{2}{\sqrt{2}}=\frac{1}{2}[/math] because [math]2^{\frac{1}{2}}=\sqrt{2}[/math]
[math]\log_{8}{4}=\frac{2}{3}[/math] because [math]8^{\frac{2}{3}}=(\sqrt{8})^2=2^2=4[/math]
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Updated Apr 21
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Mark Eichenlaub
physics curriculum developer at Art of Problem Solving
·
Upvoted by
Jay Wacker
, theoretical physicist and
James Moosh
, PhD in Pure Maths from the University of Leeds · Author has 393 answers and 10.2M answer views
·
13y
Related
What do logarithms mean in layman's terms?
tl;dr It's how many digits a number has. Since multiplying a three-digit number times a four-digit number gives a seven-digit number (and 7 = 3+4), logarithms turn multiplication into addition.
To start, it's how many zeroes come at the end of a number.
[math]\log(10) = 1[/math]
[math]\log(100) = 2[/math]
[math]\log(1000) = 3[/math]
etc.
This is useful when we want to express things in scientific notation and talk about huge numbers. If there are 10^28 atoms in a system, we can simply say that the logarithm is 28.
We also have a sense that a human is "half way" between the world of atoms (very tiny) and that of the galaxy (huge). But
tl;dr It's how many digits a number has. Since multiplying a three-digit number times a four-digit number gives a seven-digit number (and 7 = 3+4), logarithms turn multiplication into addition.
To start, it's how many zeroes come at the end of a number.
[math]\log(10) = 1[/math]
[math]\log(100) = 2[/math]
[math]\log(1000) = 3[/math]
etc.
This is useful when we want to express things in scientific notation and talk about huge numbers. If there are 10^28 atoms in a system, we can simply say that the logarithm is 28.
We also have a sense that a human is "half way" between the world of atoms (very tiny) and that of the galaxy (huge). But literally half way between them is, well, half a galaxy. If we measure the logarithm of their sizes, we find that an atom comes out to about -30, a human -1, and the galaxy around 50. A human is now much closer to being a reasonable midpoint. (I'm dividing volume by cubic meters.)
Examining our list of logarithms, we see a rule that every time we multiply by ten, we add one to the logarithm. If we divide by ten, we subtract one from the logarithm.
That rule lets us figure out log(1). It's zero, as you should check. If we keep dividing by ten, we find:
[math]\log(1) = 0[/math]
[math]\log(0.1) = -1[/math]
[math]\log(0.01) = -2[/math]
etc.
That takes care of the powers of ten, but we would like to know the values for intermediate numbers. What is log(55)?
We might say that 55 is half way between 10 and 100, so the logarithm should be half way between 1 and 2, or 1.5.
That's a decent approximation, but it's not quite right. The reason is that logarithm isn't a straight line. It increases more steeply from 1 to 10 than from 10 to 100. If we take the values we know for the powers of ten and plot them, then draw a smooth line between, we get something like this:
log(55) looks like it is more than 1.5. Perhaps 1.7.
What we might notice is that, for everything done so far,
[math]\log(10^x) = x[/math]
We can use that as our general rule for finding out logarithms. For example, 10^1.7 is 50.1 and 10^1.8 is 63.1, so log(55) should be between 1.7 and 1.8. If we carry this analysis further, we find log(55) = 1.7403...
Since there is a solution to
[math]10^x = y[/math]
for any positive number [math]y[/math], we can take the logarithm of any positive number. The logarithms themselves go from -infinity (for very, very small numbers) to infinity (for very, very big numbers). However, we can't take the logarithm of zero or a negative number. (At least not in the real number system.)
An important property of logarithms stems from the rule
[math]10^x 10^y = 10^{x+y}[/math]
that means
[math]\log(10^x 10^y) = x+y = \log(10^x) + \log(10^y)[/math]
The first and last steps are the important ones. If we represent [math]10^x[/math] as simply [math]a[/math] and [math]10^y[/math] as [math]b[/math], we can write
[math] \log(a b) = \log(a) + \log(b)[/math]
This is pretty much like saying that a three-digit number times a four-digit number is a seven-digit number. Logarithms turn multiplication into addition. This is the basic idea behind how a slide rule computes multiplication, for example.
If all this made sense, you should be able to figure out
[math]\log(\sqrt{10})[/math]
Also, if I tell you that log(2) = 0.301, you should be able to figure out the following:
[math]\log(5)[/math]
[math]\log(16)[/math]
[math]\log(1.25)[/math]
(You can check your answers with the Google search bar.)
Another, equivalent way to view logarithms is that the logarithm of a number is how many bits you need to represent it. A single bit has two states, conventionally 0 and 1. Thus, to represent two states we need a single bit, so
[math]\log_2(2) = 1[/math]
To represent four states, we can use 00, 01, 10, and 11. That's two bits, so
[math]\log_2(4) = 2[/math]
[math]\log_2(8) = 3[/math]
etc.
This is actually very similar to how we defined logarithms before, but we're using powers of two rather than powers of ten. We say the logarithm has a different base. That's why we wrote [math]\log_2[/math] rather than just [math]\log[/math]. It doesn't really matter because all the same laws of logarithms are still true.
Different bases are useful in different circumstances. Since we communicate numbers in base ten, we typically use base ten logarithms. But in information theory, base two logarithms are more natural. For limiting processes, a third base, the base e, turns out to be most natural (and is called the natural logarithm).
check out the wikipedia article:
or Napier's original work on logarithms:
or my answer on calculus and logarithms Mark Eichenlaub's answer to Why is \displaystyle\int \frac{\mathrm{d}x}{x} = \ln x ?
Job Bouwman
PhD candidate in medical imaging (MRI). Former high school math teacher.
·
Upvoted by
Jack Huizenga
, Math Professor, Penn State · Author has 1.1K answers and 7.1M answer views
·
Updated 9y
Related
What is the intuition behind the logarithm?
The logarithm counts the number of groupings.
Suppose a bakery puts 12 cookies in a package, and places 12 of these packages in a larger box for transport:
Then a box can be seen as cookies which have been grouped twice: one box contains [math]12^2[/math] = [math]144[/math] cookies. Inversely, when ordering 144 cookies, and knowing that this bakery works with base 12, the logarithm will return the number of groupings: [math]\log_{12}{144} = 2[/math] groupings.
Suppose the transport company also works in base twelve; 12 boxes are wrapped in plastic, 12 plastic units are stacked onto a wooden pallet, 12 pallets are transported in a v
The logarithm counts the number of groupings.
Suppose a bakery puts 12 cookies in a package, and places 12 of these packages in a larger box for transport:
Then a box can be seen as cookies which have been grouped twice: one box contains [math]12^2[/math] = [math]144[/math] cookies. Inversely, when ordering 144 cookies, and knowing that this bakery works with base 12, the logarithm will return the number of groupings: [math]\log_{12}{144} = 2[/math] groupings.
Suppose the transport company also works in base twelve; 12 boxes are wrapped in plastic, 12 plastic units are stacked onto a wooden pallet, 12 pallets are transported in a van, which makes [math]12^3=1728[/math] boxes per ride:
Now the total number of cookies in one transport, is obtained by multiplying the number of cookies per box with the number of boxes per ride: [math]144 \times 1728 = 248832 [/math]cookies.
However, if we look at this on a grouping scale (logarithmic scale) then we must use addition instead: the total number of groupings is [math]2 + 3 = 5[/math]. Stated otherwise: [math]\log_{12}{144} + \log_{12}{1728} = \log_{12}{248832}[/math].
Now suppose 20736 packages of cookies were stolen. Then we know that this involved [math]\log_{12}{20736} = 4[/math] groupings. So counting upwards from these packages...
...we know that this monster learned how to drive:
Nathan Pflueger
Assistant Professor of Mathematics
·
Upvoted by
David Joyce
, Ph.D. Mathematics, University of Pennsylvania (1979)
·
12y
Related
What is the intuition behind the logarithm?
The logarithm is a device that turns multiplication into addition.
Much of their early importance comes from this fact: you often want to multiply, but it turns out that adding is much easier. That's why the slide rule was such a great tool for so long.
Slide rules do addition in a simple mechanical way: if you slide a length a segment up to a length b segment, you get a length (a+b) segment. Now label the ruler with logarithms and you are doing multiplication in a simple mechanical way instead!
Abstractly, you can imagine the logarithm function as an infinite sliderule. It puts every number in
The logarithm is a device that turns multiplication into addition.
Much of their early importance comes from this fact: you often want to multiply, but it turns out that adding is much easier. That's why the slide rule was such a great tool for so long.
Slide rules do addition in a simple mechanical way: if you slide a length a segment up to a length b segment, you get a length (a+b) segment. Now label the ruler with logarithms and you are doing multiplication in a simple mechanical way instead!
Abstractly, you can imagine the logarithm function as an infinite sliderule. It puts every number in exactly the right place so that you can multiply any two by adding their horizontal positions.
Most properties of the logarithm are easy to see from this interpretation. For example: exponentiation is repeated multiplication; it becomes repeated addition. Physically: the powers of a number are evenly spaced on the sliderule.
More abstractly: logarithms exploit the fact that "positive numbers with multiplication" has exactly the same structure as "all numbers with addition." The dictionary between the two structures is the logarithm function. In the jargon, we say that the two structures are isomorphic.
Michael Lamar
PhD in Applied Mathematics
·
Upvoted by
David Joyce
, Ph.D. Mathematics, University of Pennsylvania (1979) · Author has 3.7K answers and 17.4M answer views
·
7y
Related
Why is e the base of the natural logarithm?
The defining property of the natural logarithm (as opposed to some other logarithm) is that for [math]x>0[/math]:
[math]\ln(x) = \displaystyle\int_1^x \frac {dt}t[/math]
Furthermore, for [math]a>0[/math], we define the [math]\log_a(x)=\frac{\ln (x)}{\ln( a)}[/math].
So we see that:
[math]\log_e(x)=\ln (x) \text{ if and only if }\ln (e) = 1[/math]
So your question is equivalent to, “Why is the natural log of [math]e[/math] equal to one?”
So we must prove that:
[math]\displaystyle\int_1^e \frac {dt}t=1[/math]
But it turns out that [math]e[/math] is often DEFINED to be the number for which
[math]\displaystyle\int_1^e \frac {dt}t=1[/math]
So if we choose this definition of [math]e[/math], we discover that there is nothing to prove.
The defining property of the natural logarithm (as opposed to some other logarithm) is that for [math]x>0[/math]:
[math]\ln(x) = \displaystyle\int_1^x \frac {dt}t[/math]
Furthermore, for [math]a>0[/math], we define the [math]\log_a(x)=\frac{\ln (x)}{\ln( a)}[/math].
So we see that:
[math]\log_e(x)=\ln (x) \text{ if and only if }\ln (e) = 1[/math]
So your question is equivalent to, “Why is the natural log of [math]e[/math] equal to one?”
So we must prove that:
[math]\displaystyle\int_1^e \frac {dt}t=1[/math]
But it turns out that [math]e[/math] is often DEFINED to be the number for which
[math]\displaystyle\int_1^e \frac {dt}t=1[/math]
So if we choose this definition of [math]e[/math], we discover that there is nothing to prove.
Now, if we choose some other definition for [math]e[/math] like
[math]e=\displaystyle\lim_{x\to\infty}\left (1+\frac 1x\right)^x[/math]
[math]e=\text{argmax}\left(\sqrt[x]x\right)[/math] over any [math]x>0[/math]
or
[math]e = \displaystyle\sum_{k=0}^\infty \frac 1{k!}[/math]
or any one of many other choices, we are left to show that our chosen definition of [math]e[/math] implies [math]\displaystyle\int_1^e \frac {dt}t=1[/math]. Depending on your chosen definition, that might be relatively simple to prove or it might be very challenging. But unless you tell me which definition you prefer, I don’t know what you would actually need to show to establish the truth of the claim. And since I don’t know which definition you want, I will just choose the one that is the simplest for this problem: [math]e[/math] is the real number for which [math]\displaystyle\int_1^e \frac {dt}t=1[/math].
Related questions
What is the fundamental concept behind logarithms?
What is an intuitive explanation for exponential and logarithm curves?
What is the intuition behind the logarithm?
What is the mathematical meaning of a natural logarithm?
What is a logarithmic basis?
Why is a natural logarithm called so?
What do logarithms mean in layman's terms?
What is an easy and practical way to understand or explain logarithms and logarithmic scales?
What is a logarithm (with examples)?
What is the proper explanation of the concept of logarithms?
What is the easiest way to understand logarithms?
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System of equations has no solution
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If the system of linear equations,
⎧⎩⎨x x x+a y+2 y+5 y+z+2 z+3 z=3=6=b{x+a y+z=3 x+2 y+2 z=6 x+5 y+3 z=b
has no solution, then:
a=−1,b=9 a=−1,b=9
a=−1,b≠9 a=−1,b≠9
a≠−1,b=9 a≠−1,b=9
a=1,b≠9 a=1,b≠9
I really fail to understand why the answer cannot be option (1).
Here's my method:
For a=−1 a=−1, the coefficient determinant determinant is 0. And if we create Δ z Δ z, the determinant in which coefficients of z z are replaced by 3,6,b 3,6,b, (b= 9), that determinant is non zero.
My book says that if coefficient determinant is 0 0 and at least one of Δ y,Δ z,Δ x Δ y,Δ z,Δ x is non 0 0 then the system has no solution.
But the answer given is (2)(2). What's wrong with my method? What's the best way to tackle such problems?
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edited Jan 3, 2019 at 16:33
gt6989b
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asked Jan 2, 2019 at 21:57
ArcherArcher
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For choice (1), isn't (x, y, z) = (4, 1, 0) a solution?T.J. Gaffney –T.J. Gaffney 2019-01-02 22:07:58 +00:00 Commented Jan 2, 2019 at 22:07
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You miscalculated Δ z Δ z. The book solution is correct. In order to have no solution, it is necessary that the determinant of the matrix of coefficients be zero: this indeed yields a=−1 a=−1. The rank of the matrix is 2 2 in this case, and up to a constant multiple, the only nontrivial "dependence" of the equations can be 1⋅(1)−2⋅(2)+1⋅(3)=0 1⋅(1)−2⋅(2)+1⋅(3)=0. So under the a=−1 a=−1 assumption, there are infinitely many solutions iff this generalizes to the right hand sides, that is, iff 3−2⋅6+b=0 3−2⋅6+b=0, i.e., b=9 b=9. Hence, the answer is a=−1,b≠9 a=−1,b≠9.
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answered Jan 2, 2019 at 22:16
A. PongráczA. Pongrácz
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I suppose the best way to tackle is to row-reduce the matrix and see what happens. You have to eliminate the third and the first columns:
⎛⎝⎜1 1 1 a 2 5 1 2 3 3 6 b⎞⎠⎟→⎛⎝⎜1 0 0 a 2−a 5−a 1 1 2 3 3 b−3⎞⎠⎟→⎛⎝⎜1 0 0 2 a−2 2−a 1+a 0 1 0 0 3 b−9⎞⎠⎟(1 a 1 3 1 2 2 6 1 5 3 b)→(1 a 1 3 0 2−a 1 3 0 5−a 2 b−3)→(1 2 a−2 0 0 0 2−a 1 3 0 1+a 0 b−9)
so if a=−1 a=−1 the last equation reads 0=b−9 0=b−9, thus having b=9 b=9 you have 0=0 0=0, which eliminates the equation, and now you have 2 equations in 3 unknowns, hence an infinite number of solutions. Thus (i) cannot be the answer.
The condition for no solutions would require that a=−1 a=−1 (forcing the last equation to read 0=b−9 0=b−9) and b≠9 b≠9 (forcing it to be false).
UPDATE
In retrospect, as you say, if a=−1 a=−1 then Δ=0 Δ=0, so there will not be a unique solution. However, if b=9 b=9, the last column of Δ z Δ z is (3,6,9)=3×(1,2,3)(3,6,9)=3×(1,2,3), so exactly 3 3 times the last column of Δ Δ, hence, Δ z=3 Δ=0 Δ z=3 Δ=0, and the idea in your book is correct.
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edited Jan 2, 2019 at 22:24
answered Jan 2, 2019 at 22:15
gt6989bgt6989b
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Question asks the condition for no solution...Archer –Archer 2019-01-02 22:16:10 +00:00 Commented Jan 2, 2019 at 22:16
@Abcd see update gt6989b –gt6989b 2019-01-02 22:18:08 +00:00 Commented Jan 2, 2019 at 22:18
"you have to eliminate the third and the first columns"- why?Archer –Archer 2019-01-02 22:22:58 +00:00 Commented Jan 2, 2019 at 22:22
@Abcd because they have no variables, much easier to eliminate these than pivoting on column with a a in. You can still do it, but it is more arithmetic. BTW method of your book is correct, see another update gt6989b –gt6989b 2019-01-02 22:25:24 +00:00 Commented Jan 2, 2019 at 22:25
Thanks for your time. I understand that you have applied transformations to the determinant. But I am unable to follow after that... What is the new determinant supposed to indicate?Archer –Archer 2019-01-02 22:27:47 +00:00 Commented Jan 2, 2019 at 22:27
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I would do it this way, using row reduction. The criterion for a non-homogeneous linear system to have solutions is that the matrix of the linear system (l.h.s.) and the augmented matrix have the same rank.
From this criterion we deduce readily that if this system has no solution, the matrix of the l.h.s. cannot have rank 3 3, hence its determinant, which is equal to −(a+1)−(a+1) is zero, whence a=−1 a=−1. Computing the determinant by row reduction shows this matrix has rank 2 2.
Now, if a=−1 a=−1, the augmented matrix reduces as follows:
⎡⎣⎢1 1 1−1 2 5 1 2 3 3 6 b⎤⎦⎥⇝⎡⎣⎢1 0 0−1 3 3 1 1 1 3 3 b−6⎤⎦⎥⇝⎡⎣⎢1 0 0−1 3 0 1 1 0 3 3 b−9⎤⎦⎥[1−1 1 3 1 2 2 6 1 5 3 b]⇝[1−1 1 3 0 3 1 3 0 3 1 b−6]⇝[1−1 1 3 0 3 1 3 0 0 0 b−9]
Thus, for a=−1 a=−1, the augmented matrix has rank 3 3 if and only if b−9≠0 b−9≠0.
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answered Jan 2, 2019 at 22:24
BernardBernard
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Rewrite the system as:
A X=B A X=B
where X=(x,y,z)t X=(x,y,z)t, B=(3,6,b)t B=(3,6,b)t and
A=⎡⎣⎢1 1 1 a 2 5 1 2 3⎤⎦⎥A=[1 a 1 1 2 2 1 5 3]
First you must have d e t(A)=0 d e t(A)=0 which after calculations give a=−1 a=−1. So the system becomes:
A=⎡⎣⎢1 1 1−1 2 5 1 2 3⎤⎦⎥A=[1−1 1 1 2 2 1 5 3]
Now, row 3 becomes all zero when r 3→r 3−2 r 2+r 1 r 3→r 3−2 r 2+r 1. Applying this to the augmented matrix [A|B][A|B] gives the (3,4)(3,4) entry to be:
b−9 b−9
It is required that this entry is nonzero in order to have an inconsistent matrix. Thus, a=−1,b≠9 a=−1,b≠9.
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answered Jan 2, 2019 at 22:27
Kal S.Kal S.
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13758 | https://artofproblemsolving.com/wiki/index.php/Modular_arithmetic/Introduction?srsltid=AfmBOoqYnhLAh9s6Ha7dg6RkR4HAZkzD-jcg3PWJtzyvHsVwwaF0Dwpi | Art of Problem Solving
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Modular arithmetic/Introduction
Modular arithmetic is a special type of arithmetic that involves only integers. This goal of this article is to explain the basics of modular arithmetic while presenting a progression of more difficult and more interesting problems that are easily solved using modular arithmetic.
Contents
[hide]
1 Introductory Video
2 Understand Modular Arithmetic
3 Residue
4 Congruence
4.1 Examples
4.2 Sample Problem
4.2.1 Solution:
4.2.2 Another Solution:
4.2.3 Another Solution:
5 Making Computation Easier
5.1 Addition
5.1.1 Problem
5.1.2 Solution
5.1.3 Why we only need to use remainders
5.1.4 Solution using modular arithmetic
5.1.5 Addition rule
5.1.6 Proof of the addition rule
5.2 Subtraction
5.2.1 Problem
5.2.2 Solution
5.2.3 Subtraction rule
5.3 Multiplication
5.3.1 Problem
5.3.2 Solution
5.3.3 Solution using modular arithmetic
5.3.4 Multiplication rule
5.4 Exponentiation
5.4.1 Problem #1
5.4.2 Problem #2
5.4.3 Problem #3
6 Summary of Useful Facts
7 Problem Applications
8 Applications of Modular Arithmetic
9 Resources
10 See also
Introductory Video
Understand Modular Arithmetic
Let's use a clock as an example, except let's replace the at the top of the clock with a .
This is the way in which we count in modulo 12. When we add to , we arrive back at . The same is true in any other modulus (modular arithmetic system). In modulo , we count
We can also count backwards in modulo 5. Any time we subtract 1 from 0, we get 4. So, the integers from to , when written in modulo 5, are
where is the same as in modulo 5. Because all integers can be expressed as , , , , or in modulo 5, we give these integers their own name: the residue classes modulo 5. In general, for a natural number that is greater than 1, the modulo residues are the integers that are whole numbers less than :
This just relates each integer to its remainder from the Division Theorem. While this may not seem all that useful at first, counting in this way can help us solve an enormous array of number theory problems much more easily!
Residue
We say that is the modulo-residue of when , and .
Congruence
There is a mathematical way of saying that all of the integers are the same as one of the modulo 5 residues. For instance, we say that 7 and 2 are congruent modulo 5. We write this using the symbol : In other words, this means in base 5, these integers have the same residue modulo 5:
The (mod 5) part just tells us that we are working with the integers modulo 5. In modulo 5, two integers are congruent when their difference is a multiple of 5. In general, two integers and are congruent modulo when is a multiple of . In other words, when is an integer. Otherwise, , which means that and are not congruent modulo .
Examples
because is a multiple of .
because , which is an integer.
because , which is not a multiple of .
because , which is not an integer.
Sample Problem
Find the modulo residue of .
Solution:
Since R , we know that
and is the modulo residue of .
Another Solution:
Since , we know that
We can now solve it easily
and is the modulo residue of
Another Solution:
We know is a multiple of since is a multiple of . Thus, and is the modulo residue of .
Making Computation Easier
We don't always need to perform tedious computations to discover solutions to interesting problems. If all we need to know about are remainders when integers are divided by , then we can work directly with those remainders in modulo . This can be more easily understood with a few examples.
Addition
Problem
Suppose we want to find the units digit of the following sum:
We could find their sum, which is , and note that the units digit is . However, we could find the units digit with far less calculation.
Solution
We can simply add the units digits of the addends:
The units digit of this sum is , which must be the same as the units digit of the four-digit sum we computed earlier.
Why we only need to use remainders
We can rewrite each of the integers in terms of multiples of and remainders:
.
When we add all four integers, we get
At this point, we already see the units digits grouped apart and added to a multiple of (which will not affect the units digit of the sum):
.
Solution using modular arithmetic
Now let's look back at this solution, using modular arithmetic from the start. Note that
Because we only need the modulo residue of the sum, we add just the residues of the summands:
so the units digit of the sum is just .
Addition rule
In general, when , and are integers and is a positive integer such that
the following is always true:
.
And as we did in the problem above, we can apply more pairs of equivalent integers to both sides, just repeating this simple principle.
Proof of the addition rule
Let , and where and are integers. Adding the two equations we get:
Which is equivalent to saying
Subtraction
The same shortcut that works with addition of remainders works also with subtraction.
Problem
Find the remainder when the difference between and is divided by .
Solution
Note that and . So,
Thus,
so 1 is the remainder when the difference is divided by . (Perform the subtraction yourself, divide by , and see!)
Subtraction rule
When , and are integers and is a positive integer such that
the following is always true:
Multiplication
Modular arithmetic provides an even larger advantage when multiplying than when adding or subtracting. Let's take a look at a problem that demonstrates the point.
Problem
Jerry has boxes of soda in his truck. The cans of soda in each box are packed oddly so that there are cans of soda in each box. Jerry plans to pack the sodas into cases of cans to sell. After making as many complete cases as possible, how many sodas will Jerry have leftover?
Solution
First, we note that this word problem is asking us to find the remainder when the product is divided by .
Now, we can write each and in terms of multiples of and remainders:
This gives us a nice way to view their product:
Using FOIL, we get that this equals
We can already see that each part of the product is a multiple of , except the product of the remainders when each and are divided by 12. That part of the product is , which leaves a remainder of when divided by . So, Jerry has sodas leftover after making as many cases of as possible.
Solution using modular arithmetic
First, we note that
Thus,
meaning there are sodas leftover. Yeah, that was much easier.
Multiplication rule
When , and are integers and is a positive integer such that
The following is always true:
.
Exponentiation
Since exponentiation is just repeated multiplication, it makes sense that modular arithmetic would make many problems involving exponents easier. In fact, the advantage in computation is even larger and we explore it a great deal more in the intermediate modular arithmetic article.
Note to everybody: Exponentiation is very useful as in the following problem:
Problem #1
What is the last digit of if there are 1000 7s as exponents and only one 7 in the middle?
We can solve this problem using mods. This can also be stated as . After that, we see that 7 is congruent to -1 in mod 4, so we can use this fact to replace the 7s with -1s, because 7 has a pattern of repetitive period 4 for the units digit. is simply 1, so therefore , which really is the last digit.
Problem #2
What are the tens and units digits of ?
We could (in theory) solve this problem by trying to compute , but this would be extremely time-consuming. Moreover, it would give us much more information than we need. Since we want only the tens and units digits of the number in question, it suffices to find the remainder when the number is divided by . In other words, all of the information we need can be found using arithmetic mod .
We begin by writing down the first few powers of mod :
A pattern emerges! We see that So for any positive integer , we have (mod ). In particular, we can write
.
By the "multiplication" property above, then, it follows that
(mod ).
Therefore, by the definition of congruence, differs from by a multiple of . Since both integers are positive, this means that they share the same tens and units digits. Those digits are and , respectively.
Problem #3
Can you find a number that is both a multiple of but not a multiple of and a perfect square?
No, you cannot. Rewriting the question, we see that it asks us to find an integer that satisfies .
Taking mod on both sides, we find that . Now, all we are missing is proof that no matter what is, will never be a multiple of plus , so we work with cases:
This assures us that it is impossible to find such a number.
Summary of Useful Facts
Consider four integers and a positive integer such that and . In modular arithmetic, the following identities hold:
Addition: .
Subtraction: .
Multiplication: .
Division: , where is a positive integer that divides and .
Exponentiation: where is a positive integer.
Problem Applications
Applications of Modular Arithmetic
Modular arithmetic is an extremely flexible problem solving tool. The following topics are just a few applications and extensions of its use:
Divisibility rules
Linear congruences
Resources
The AoPS Introduction to Number Theory by Mathew Crawford.
The AoPS Introduction to Number Theory Course. Thousands of students have learned more about modular arithmetic and problem solving from this 12 week class.
See also
Intermediate modular arithmetic
Olympiad modular arithmetic
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13759 | https://www.gauthmath.com/solution/1704045096045574/10-The-roots-of-the-quadratic-equation-3x2-2x-1-0-are-and-i-State-the-values-of- | Solved: The roots of the quadratic equation 3x^2-2x+1=0 are α and β. (i) State the values of alpha [Math]
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The roots of the quadratic equation 3x^2-2x+1=0 are α and β. (i) State the values of alpha +beta and alpha beta Hence find the value of alpha^2+beta^2 [B] (ii) Show that alpha^4+beta^4=- 14/81 (iii) Find the quadratic equations whose roots are 1/alpha 4 and 1/beta 4
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eet_3_S2_2024.pdf b If a quadratic equation has roots alpha =2+ square root of 3 and beta =2- square root of 3 , find the equation. 3 c The equation 2x2-x+3=0 has roots α and β. Without solving the equation, i write down the value of alpha +beta and the value of αβ. 1 ii find the value of frac 1alpha +frac 1beta and that of frac 1beta -frac 1alpha . 4 iii find a quadratic equation which has roots 2alpha -frac 1beta and 2beta -frac 1alpha . 4 d The equation x2-2x+3=0 has roots α and β. Without solving the equation, i write down the value of alpha +beta and the value of αβ. 1 ii show that alpha 2+beta 2=-2. 2 iii find the value of alpha 3+beta 3 2 iv show that alpha 4+beta 4=alpha 2+beta 22-2alpha beta 2. 2 v find a quadratic equation which has roots alpha 3-beta and beta 3-alpha , , giving your answer in the form px2+qx+r=0 where p, q and r are integers. 4 e The equation to search 20°C Partly cloudy 3:53 PM 11/10/2024
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If a quadratic equation has roots alpha =2+ square root of 3 and beta =2- square root of 3 , find the equation c The equation 2x2-x+3=0 has roots α and β. Without solving the equation, i write down the value of alpha +beta and the value of αβ. ii find the value of frac 1alpha +frac 1beta and that of frac 1beta -frac 1alpha . iii find a quadratic equation which has roots 2alpha -frac 1beta and 2beta -frac 1alpha . d The equation x2-2x+3=0 has roots α and β. Without solving the equation, i write down the value of alpha +beta and the value of alpha beta . ii show that alpha 2+beta 2=-2. iii find the value of alpha 3+beta 3. iv show that alpha 4+beta 4=alpha 2+beta 22-2alpha beta 2. v find a quadratic equation which has roots alpha 3-beta and beta 3-alpha , giving your answer in the form px2+qx+r=0 where p, q and r are integers.
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Tangent, OR
Sine Rule
Trigonometric Formulas
Cosine (math function)
Mathematical Functions
Tangent Function
Trigonometric Analysis
Sine and Cosine
5
Why and how exactly do sine, cosine and tangent work? What are they?
Leo Cutter
Studied Mathematics at Bishop Gorman High School (Graduated 2020)
·
Upvoted by
Terry Moore
, M.Sc. Mathematics, University of Southampton (1968) and
Yusuf Dadkhah
, master degree Mathematics (2024) · Author has 356 answers and 4.1M answer views
·
6y
Since, cosine and tangent all have many definitions but it seems that from your question you are just asking for the “rundown”.
At their cores, sine and cosine are defined to be the y and x coordinate respectively of an angle measure on the unit circle. By the way, the unit circle is defined to be the graph of the solutions to the equation x2+y2=1
So we see that if we go some angle, θ counterclockwise with respect to the positive x-axis, our radius will land on some point of the circle. The y-coordinate of where this radius lies on the circle is called the sine of the angle, or sinθ
Since, cosine and tangent all have many definitions but it seems that from your question you are just asking for the “rundown”.
At their cores, sine and cosine are defined to be the y and x coordinate respectively of an angle measure on the unit circle. By the way, the unit circle is defined to be the graph of the solutions to the equation x2+y2=1
So we see that if we go some angle, θ counterclockwise with respect to the positive x-axis, our radius will land on some point of the circle. The y-coordinate of where this radius lies on the circle is called the sine of the angle, or sinθ. The x-coordinate of where this radius lies on the circle will be called the cosine of the angle, or cosθ. When you calculate sin(30∘), all you’re doing is finding the y-coordinate of wherever the radius lands on the circumference when rotated counterclockwise 30∘.
From here, we can define tangent. Usually most textbooks and highschools just define this to be
but this isn’t satisfying. Sure, it’s correct, but it doesn’t really explain what that has to do with anything. After all, isn’t a tangent line a line that touches a graph once? Why does the trigonometric function tangent called what it is if the two seem to have nothing to do with each other?
So this picture outlines what the tangent is; the length of a part of a line segment that is tangent to the unit circle at that given degree measure. When you’re calculating , all you’re doing is finding the length of this part of the line segment tangent to the circle when the point of tangency is located up from the positive -axis.
Here is secant of the angle represented pictorially on the unit circle.
And that is everything you could ever want when representing the 6 trigonometric functions on the unit circle.
Buddha Buck
Took calculus as an undergraduate
·
Upvoted by
Terry Moore
, M.Sc. Mathematics, University of Southampton (1968) · Author has 5.8K answers and 16.9M answer views
·
9y
Originally Answered: What is the relationship between sine, cosine, and tangent, explained simply?
·
For the purposes of my discussion, there are three main trigonometric functions: sine, tangent, and secant. I know that you didn't ask about secant, and I'm leaving out cosine, but bear with me, I'll get back to them.
Draw a circle with a radius of 1, centered at the origin, and draw a line from the origin through the first quadrant (where both and are positive). Where the line crosses the circle, draw another line perpendicular to the first. Finally, draw a line segment from where the line crosses the circle to the x axis, and from where the line crosses the circle to the y axis.
There are
For the purposes of my discussion, there are three main trigonometric functions: sine, tangent, and secant. I know that you didn't ask about secant, and I'm leaving out cosine, but bear with me, I'll get back to them.
Draw a circle with a radius of 1, centered at the origin, and draw a line from the origin through the first quadrant (where both and are positive). Where the line crosses the circle, draw another line perpendicular to the first. Finally, draw a line segment from where the line crosses the circle to the x axis, and from where the line crosses the circle to the y axis.
There are several similar right triangles here. The blue triangle (bounded by the radius, the segment marked "tangent", and the segment marked "secant"), is similar to the pink triangle (bounded by the radius, the segment marked "cosecant", and the segment marked "cotangent"). Also similar is the dark blue triangle (bounded by "sine", "cosine", and the radius).
The line with the labels "tangent" and "cotangent" forms a tangent line to the unit circle -- a line which shares just a single point with the circle. This is why the trigonometric identity associated with it is called the tangent.
"Secant" comes from the Greek meaning to cut through, and the secant line cuts through the circle -- it doesn't stop at the circle, and goes from inside to outside.
"Sine" comes from a series of mistranslations, but ultimately comes from a sanskrit word meaning bowstring.
The "co-" prefix means "complementary", and refers to the complementary angle to the measured angle. The angle of the pink corner next to the origin in the picture is the complementary angle to . You can see that the line labeled "cosecant" is the secant of the pink corner, and the line labeled "cotangent" is the tangent of the pink corner. With a little effort, it should be easy to see that the line labeled "cosine" is the sine of the pink corner as well. (There are also obsolete or obscure trig functions like the "versed sine" (versine), which also have "co-" versions, like the coversine. But since you can "verse" the cosine as well, you also have the vercosine and the covercosine. I have no idea what they are good for. Half the versed sine (the "haversine") is used in some traditional nautical navigation formulas. I don't know if anyone ever used the cohavercosine.)
The nice thing about this diagram is that it also captures all the major trig identities, using similar triangles and the pythagorean theorem. Since radius-sin-cosine form a right triangle, you get . Since the pink and blue triangles are similar and share a side of length 1, you can use that to show that the cotangent is the reciprocal of the tangent, the cosecant is the reciprocal of the sine, etc.
Matt Hodel
Former Theoretical Physics Graduate Student at McGill University (2017–2019)
·
Upvoted by
David Joyce
, Ph.D. Mathematics, University of Pennsylvania (1979) · Author has 257 answers and 1M answer views
·
11y
Originally Answered: What are the sine, cosine, and tangent functions?
·
As Paulo touches upon, the simplest way to understand the sine and cosine functions is to imagine a point moving around the unit circle.
Draw a radial line from the center of the circle to the point, as shown in the picture. Call the angle this line makes with the x-axis , and let the coordinates of the point (x,y). Then the definitions of the sine and cosine functions are
The tangent function can be defined as the ratio of the sine to the cosine:
It is also equal to the slope of that radial line you drew. It is also equal to the
As Paulo touches upon, the simplest way to understand the sine and cosine functions is to imagine a point moving around the unit circle.
Draw a radial line from the center of the circle to the point, as shown in the picture. Call the angle this line makes with the x-axis , and let the coordinates of the point (x,y). Then the definitions of the sine and cosine functions are
The tangent function can be defined as the ratio of the sine to the cosine:
It is also equal to the slope of that radial line you drew. It is also equal to the length of the blue line in the figure (can you see why?)
What you learn in school, that cosine and sine are the ratios of the opposite and adjacent legs to the hypotenuse of a right triangle, respectively, should be clear from the figure (here the hypotenuse is conveniently equal to one).
Related questions
What is the formula for cosine, sine, and tangent?
What actually are sine, cosine and tangent?
How do I know when to use ‘sine’, ‘cosine’, or ‘tangent’?
What are the uses of sine, cosine, and tangent in real life?
What do sine, cosine, and tangent actually do?
Assistant
Bot
·
1y
Sine, cosine, and tangent are fundamental functions in trigonometry that relate the angles and sides of right triangles. They are defined based on the ratios of the sides of a right triangle, but they also have broader applications in various fields, including physics, engineering, and computer science.
Definitions
Sine (sin): In a right triangle, the sine of an angle is the ratio of the length of the side opposite the angle to the length of the hypotenuse (the longest side of the triangle).
Cosine (cos): The cosine of an an
Sine, cosine, and tangent are fundamental functions in trigonometry that relate the angles and sides of right triangles. They are defined based on the ratios of the sides of a right triangle, but they also have broader applications in various fields, including physics, engineering, and computer science.
Definitions
Sine (sin): In a right triangle, the sine of an angle is the ratio of the length of the side opposite the angle to the length of the hypotenuse (the longest side of the triangle).
Cosine (cos): The cosine of an angle is the ratio of the length of the adjacent side (the side next to the angle) to the length of the hypotenuse.
Tangent (tan): The tangent of an angle is the ratio of the length of the side opposite the angle to the length of the adjacent side.
Unit Circle Definition
Beyond right triangles, sine, cosine, and tangent can also be defined using the unit circle, which is a circle with a radius of 1 centered at the origin of a coordinate system.
Sine: For an angle measured from the positive x-axis, the sine of is the y-coordinate of the point where the terminal side of the angle intersects the unit circle.
Cosine: The cosine of is the x-coordinate of the same point on the unit circle.
Tangent: Since tangent is the ratio of sine to cosine, it can also be interpreted as the slope of the line formed by the terminal side of the angle, or as the length of the segment from the origin to the point where the line intersects the tangent line at .
Properties
Periodicity: Sine and cosine functions are periodic with a period of (or 360 degrees), meaning they repeat their values in cycles. Tangent has a period of (or 180 degrees).
Range:
Sine and cosine values range from -1 to 1.
Tangent values can take any real number from to .
Pythagorean Identity: A fundamental relationship among these functions is given by the Pythagorean identity:
This identity reflects the relationship between the lengths of the sides of a right triangle.
Applications
Sine, cosine, and tangent are used in various applications, including:
Wave Motion: These functions model periodic phenomena like sound waves, light waves, and alternating current.
Physics: They are used to analyze forces, motion, and oscillations.
Engineering: Trigonometric functions are crucial in design and analysis of structures, circuits, and mechanical systems.
Computer Graphics: They help in rendering shapes, animations, and rotations.
In summary, sine, cosine, and tangent are essential mathematical functions derived from the geometry of triangles and the unit circle, with extensive applications across science and engineering.
RanjeetKumar Vishwakarma
M.Sc. in Engineering Mathematics, Institute of Chemical Technology, Mumbai (Graduated 2019)
·
Upvoted by
Terry Moore
, M.Sc. Mathematics, University of Southampton (1968)
·
8y
Originally Answered: What are the sine, cosine, and tangent functions in trigonometry?
·
The trigonometric functions (also called the circular functions) are function of an angle. They relate the angles of a triangle to the lengths of its sides. Trigonometric functions are important in the study of triangles and modeling periodic phenomena, among many other applications.
This function we can observe by Unit circle.
Animation showing how the sine function (in red) is graphed from the y-coordinate (red dot) of a point on the unit circle (in green) at an angle of
θ
in radians .
A unit circle is the circle of radius one centered at the origin (0, 0) in the cartesian coordinate system.
Le
The trigonometric functions (also called the circular functions) are function of an angle. They relate the angles of a triangle to the lengths of its sides. Trigonometric functions are important in the study of triangles and modeling periodic phenomena, among many other applications.
This function we can observe by Unit circle.
Animation showing how the sine function (in red) is graphed from the y-coordinate (red dot) of a point on the unit circle (in green) at an angle of θ in radians .
A unit circle is the circle of radius one centered at the origin (0, 0) in the cartesian coordinate system.
Let a line through the origin, making an angle of θ with the positive half of the x-axis, intersect the unit circle. The x- and y-coordinates of this point of intersection are equal to cos θ and sin(θ), respectively. The point's distance from the origin is always 1.
Unlike the definitions with the right triangle or slope, the angle can be extended to the full set of real arguments by using the unit circle. This can also be achieved by requiring certain symmetries and that sine be a periodic function.
Here,we use right angle triangle for determining the trigonometric function.
The most familiar trigonometric functions are the sine, cosine, and tangent.
Sine, Cosine and Tangent are the main functions used in trigonometry and are based on a Right angle triangle.
Here We can see name of the each side of triangle with angle θ.
· "Opposite" is opposite to the angle θ
· "Adjacent" is adjacent (next to) to the angle θ
· "Hypotenuse" is the long one
Here We can observe by another diagram.
Sine, Cosine and Tangent (often shortened to sin, cos and tan) are each a ratio of sides of a right angled triangle:
These are the some Ideas about trigonometric function.
hope you understand…
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Mervyn Tong
Studied at University of Cambridge (Graduated 2021)
·
9y
Originally Answered: What is the relationship between sine, cosine, and tangent, explained simply?
·
To get a general idea of sine, cosine and tangent, we must depart from the right-angled triangle and look at the unit circle instead. Consider this figure (credits to Stack Exchange):
Any point on this unit circle can be represented by the coordinates
math[/math]
, where
[math]\theta[/math]
is the angle anti-clockwise from the positive side of the x-axis to the radius connecting the point to the centre.
[math]tan\theta[/math]
would be the slope of the line, as
[math]tan\theta = \frac{sin\theta}{cos\theta}[/math]
by definition.
To get a general idea of sine, cosine and tangent, we must depart from the right-angled triangle and look at the unit circle instead. Consider this figure (credits to Stack Exchange):
Any point on this unit circle can be represented by the coordinates
math[/math]
, where
[math]\theta[/math]
is the angle anti-clockwise from the positive side of the x-axis to the radius connecting the point to the centre.
[math]tan\theta[/math]
would be the slope of the line, as
[math]tan\theta = \frac{sin\theta}{cos\theta}[/math]
by definition.
Related questions
Who first used the terms sine, cosine, and tangent?
How are sine, cosine, tangent, and secant related to each other (mathematically)?
What is the sine, cosine, and tangent of 30 degrees?
What is the practical meaning and use of a cosine, sine, and tangent in mathematics?
How did students find sine cosine and tangent before calculators?
Philip Lloyd
Specialist Calculus Teacher, Motivator and Baroque Trumpet Soloist.
·
Author has 6.8K answers and 52.8M answer views
·
5y
Originally Answered: What is the logic behind sin, cos, and tangent exactly?
·
I think this is the logic you are requiring.
Basically, trigonometry of right angled triangles is just a special case of “similar triangles”.
Let me explain…
Two triangles are similar if one is an enlargement of the other!
(The angles are equal but the sides are in Proportion)
The triangles need not have right angles but in this case they do to make the working obvious and the link to trigonometry clearer too.
(In fact we say that the scale factor is 2)
The working out is extremely elementary…
Now consider these similar triangles…
Long before we had calculators, people made tables for lots of values of
I think this is the logic you are requiring.
Basically, trigonometry of right angled triangles is just a special case of “similar triangles”.
Let me explain…
Two triangles are similar if one is an enlargement of the other!
(The angles are equal but the sides are in Proportion)
The triangles need not have right angles but in this case they do to make the working obvious and the link to trigonometry clearer too.
(In fact we say that the scale factor is 2)
The working out is extremely elementary…
Now consider these similar triangles…
Long before we had calculators, people made tables for lots of values of sin (θ) and cos (θ) for all the angles from 0 to 90 degrees.
If we had a triangle with a Hypotenuse of 7 cm and an angle of 53 degrees we would find the other two sides as follows…
I believe this “enlargement method” based on the “unit triangle” is more intuitive than the SOH-CAH-TOA idea but it is basically the same.
In fact tangent problems would be done the same way.
NB I think it is sad that people have to rely on formulas and mnemonics like SOH-CAH-TOA which basically do the thinking for you instead of using logic.
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Kathryn Stewart
Secondary Math Teacher
·
Upvoted by
Terry Moore
, M.Sc. Mathematics, University of Southampton (1968)
·
6y
Originally Answered: What do sine, cosine, and tangent actually do?
·
The short answer is that they succinctly name the relationships they represent in right angled triangles.
Here’s a simple way of thinking about it: you are going to lean a ladder against a wall. The sine ratio of the angle that the ladder makes with the ground compares how far up the wall your ladder reaches as a fraction of the length of the ladder. The cosine ratio is how far from the wall you’re putting the feet of the ladder compared to the length of the ladder. The tan ratio compares how far up to how far back so tells the slope of the ladder.
For example, if you put a 12 foot ladder so tha
The short answer is that they succinctly name the relationships they represent in right angled triangles.
Here’s a simple way of thinking about it: you are going to lean a ladder against a wall. The sine ratio of the angle that the ladder makes with the ground compares how far up the wall your ladder reaches as a fraction of the length of the ladder. The cosine ratio is how far from the wall you’re putting the feet of the ladder compared to the length of the ladder. The tan ratio compares how far up to how far back so tells the slope of the ladder.
For example, if you put a 12 foot ladder so that the top of it touches 6 feet above the ground, the sine ratio is 1/2 and so the angle that you will have made with the ground is 30 degrees. Your ladder probably will collapse when you step on it. Get a shorter ladder or lean it higher up the wall.
I generally put my 6 foot ladder about 2 feet from the wall. That means my cosine ratio would be 2/6 or 1/3. The angle my ladder made with the ground can be found using inverse cosine. I’m sure that someone in the land of health and safety somewhere will have determined safe angles for ladders. Since people don’t walk around with protractors, they likely communicated those angles as the range of distances from the wall you can place the feet of the ladder which would depend in part on the strength of the material the ladder is made from, and partly from its length. Regardless, in telling you how far out from the wall to put a specific ladder, they’re describing the adjacent side in the cosine ratio.
If I built a ramp that had one ending the ground 4 feet from my door that is 3 feet off the ground, my ramp would be 5m long (use Pythagoras to check). The tan ratio is the slope of my ramp (4/3). That’s quite steep! Wouldn’t be ‘to code’ for wheelchair use! I think anything more than a tan ratio of 12% might be the upper limit to use if you’re trying to make your front door wheelchair accessible.
Sine, cosine and tangent are useful in situations such as ramps, roofs and roads where there are guidelines that need to be followed.
Hope that helps put a practical spin on the primary trig ratios!
Steven Lehar
Independent researcher with radical ideas
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Author has 417 answers and 688.3K answer views
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8y
Originally Answered: What actually are sine, cosine and tangent?
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I really think math teachers should make their students generate their own trig tables before using their calculators. It is really easy!
Draw a large circle with a compass, and draw horizontal/vertical x and y axes through the center. Now for each degree from 0 to 360, draw a line from the origin to the circle at that angle, then record the x coordinate (drop a perpendicular to the x axis) and the y coordinate (drop a horizontal over to the y axis) and you have the cosine(of that angle) and the sin(of that angle). That is what they really are! (And for tangent record the ratio of the y/x for e
I really think math teachers should make their students generate their own trig tables before using their calculators. It is really easy!
Draw a large circle with a compass, and draw horizontal/vertical x and y axes through the center. Now for each degree from 0 to 360, draw a line from the origin to the circle at that angle, then record the x coordinate (drop a perpendicular to the x axis) and the y coordinate (drop a horizontal over to the y axis) and you have the cosine(of that angle) and the sin(of that angle). That is what they really are! (And for tangent record the ratio of the y/x for each angle)
More generally (this is where it gets interesting) trigonometry defines a mapping from the city-block Cartesian grid-like coordinate system to a dart-board polar angular and radial coordinate system. Trigonometry translates back and forth between the two systems.
The definition of the irrational i, the square root of minus one, is the equation that relates the grid and dart-board systems to each other by declaring that if you rotate your angle to 180 degrees (reverse direction completely in the dartboard coordinate system) that is equivalent to multiplying by -1, turning 180 degrees is equivalent to reversing your direction (negation).
I don’t care what you say, that there is awesome!
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Raelene Maths
Story Teller | Creator | Upcycler
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Author has 63 answers and 181.6K answer views
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6y
Originally Answered: What do sine, cosine, and tangent actually do?
·
Sine, cosine and tangent each compare the length of one side of a right triangle to another length. They describe the length of one side relative to another side.
On their own, sine, cosine and tangent cannot tell the difference between the following two similar right triangles:
To do that, just one length of the triangle must be given. Then the rest of the lengths can be scaled to fit. For example [math]\frac{1}{2} = \frac{3}{6}[/math] but [math]1 \neq 3[/math] and 2[math] \neq 6[/math].
For three lengths forming a right triangle inclined by 30 degrees,
[math]\sin{30^o} = {0.5}[/math] says that the height is 50% of the hypotenuse
[math]\cos{30^o} \approx[/math]
Sine, cosine and tangent each compare the length of one side of a right triangle to another length. They describe the length of one side relative to another side.
On their own, sine, cosine and tangent cannot tell the difference between the following two similar right triangles:
To do that, just one length of the triangle must be given. Then the rest of the lengths can be scaled to fit. For example [math]\frac{1}{2} = \frac{3}{6}[/math] but [math]1 \neq 3[/math] and 2[math] \neq 6[/math].
For three lengths forming a right triangle inclined by 30 degrees,
[math]\sin{30^o} = {0.5}[/math] says that the height is 50% of the hypotenuse
[math]\cos{30^o} \approx {0.8660}[/math] says that the base is 86.6% of the hypotenuse
[math]\tan{30^o} \approx {0.5774}[/math] says that the height is 57.74% of the base
Knowing that [math]\sin{30^o} = {0.5}[/math] narrows the search for the correct triangle to a family of similar right triangles, all with the same 3 angles. Then, knowing one more side uniquely identifies one triangle out of the infinite family. Trigonometry at its heart is about proportions and equivalent fractions.
Truth is, you really only need ONE of sine, cosine or tangent. But using the others makes naming everything else (in a right triangle) easier. It would be like knowing your name and then naming everyone you know by their association to you. It gets a bit cumbersome after a while.
For example, we could choose to use only [math]\sin\theta[/math] and then combine it with the Pythagorean Theorem to get [math]\sqrt{1 - \sin^2{\theta}}[/math] for [math]\cos\theta[/math] and [math]\frac{\sin\theta}{\sqrt{1 - \sin^2{\theta}}}[/math] for [math]\tan\theta[/math].
Scott Brickner
I regret not majoring in math, so I do it as a hobby
·
Upvoted by
Terry Moore
, M.Sc. Mathematics, University of Southampton (1968) · Author has 19K answers and 6.5M answer views
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6y
Originally Answered: Why do sine, cosine, and tangent work; where did these rules(?) Originate from. I use the law of cosine daily at work, and use all three functions daily, yet I have no idea what it really is.
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If you remember much of high school geometry, you’ll probably remember problems where you had to figure out the lengths or angles in various plane figures. If you were any good at solving those problems, you’d probably remember that most of the time, the answers came down to finding certain triangles where you knew some of the lengths and angles and the value you were seeking was one of the missing ones.
Trigonometry (literally “the measurement of triangles”) largely comes from that. You’ve got a triangle, which has six “parts”—three sides and three angles. You know some of them and you want to
If you remember much of high school geometry, you’ll probably remember problems where you had to figure out the lengths or angles in various plane figures. If you were any good at solving those problems, you’d probably remember that most of the time, the answers came down to finding certain triangles where you knew some of the lengths and angles and the value you were seeking was one of the missing ones.
Trigonometry (literally “the measurement of triangles”) largely comes from that. You’ve got a triangle, which has six “parts”—three sides and three angles. You know some of them and you want to know at least one of the others.
You can always make a line perpendicular to one of the sides through the opposite vertex, so if the triangle isn’t a right triangle, you can create a similar problem with two right triangles, solve that and use the answers to solve the original one. This is a useful strategy because you always know one of the angles in a right triangle.
So the core of trigonometry comes down to dealing with right triangles. If you know just one other angle of a right triangle, then the rest of the triangle is completely determined up to similarity (the ratios of the sides). This means that (up to similarity) a right triangle is uniquely identified by one of the base angles.
The strategy trigonometry takes is to just give names to the rest of the parts. You’ve got an angle, its complement, and the right angle. The sides are called the “opposite”, “adjacent” and the “hypotenuse”. The ratios of their lengths are called the “sine” (ratio of the opposite to the hypotenuse), the tangent (ratio of the opposite to the adjacent), and the secant (ratio of the hypotenuse to the adjacent).
That names the three ratios from the perspective of a given angle—but if you take the other angle (the complement of the first), then the ratios become the cosine, cosecant, and cotangent.
That’s where they come from. They’re just names given to ratios in a right triangle that’s been identified by one of the (non-right) angles.
The rest of trigonometry is showing how to use them to solve various problems. The law of cosines, for example, tells you how to work backward to your original (non-right) triangle when you got a result by dropping a perpendicular from one of the vertices.
Will Gates
Former Assistant chief cook and bottle washer.
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Author has 2K answers and 1.8M answer views
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6y
Imagine (or draw) a circle, it can have any radius: r, but, for simplicity, start with r = 1m.
Let's start with the well known Pythagorean triple: (3, 4, 5) where the 5 represents the hypotenuse and the other two sides are 3 and 4.
On our circle, the radius can be used as the hypotenuse of the triangle, hence 5 units = 1 m = 1000 mm and one unit = 200 mm. Draw a horizontal diameter and measure 3 units = 3 x 200 mm = 600 mm along the horizontal from the centre towards the right. Construct a vertical line from this point, up to the circumference of the circle.
Measure this vertical. It will be 4 un
Imagine (or draw) a circle, it can have any radius: r, but, for simplicity, start with r = 1m.
Let's start with the well known Pythagorean triple: (3, 4, 5) where the 5 represents the hypotenuse and the other two sides are 3 and 4.
On our circle, the radius can be used as the hypotenuse of the triangle, hence 5 units = 1 m = 1000 mm and one unit = 200 mm. Draw a horizontal diameter and measure 3 units = 3 x 200 mm = 600 mm along the horizontal from the centre towards the right. Construct a vertical line from this point, up to the circumference of the circle.
Measure this vertical. It will be 4 units or 800 mm. Join the top of the vertical to the centre of the circle, completing the 3, 4, 5 triangle. Label the sides a (3 units), b (4 units) and the hypotenuse c (5 units).
Now, take a photograph of the circle (using an old, film camera). When the negative is developed, it will be a scaled down image of the original drawing. The fact that it is scaled down implies the fact that the small triangle on the negative is in direct proportion to the large one on the drawing.
Measure the angles of the triangle: they will be (in degrees) 53.13010235…, 36.86989765…, 90.
(Well, ok, you might need to use trigonometry to get this accuracy. With a protractor you will probably get approximate values of: 37, 53, 90 degrees.)
If you measure the angles of the triangle in the photograph, they will be exactly the same. The triangles are similar.
Mark the photograph with apostrophes making the sides a', b' and c' (Read as "a prime", …)
These are known as corresponding sides (to a, b and c).
Measure a', b' and c',
a / a'
b / b'
c / c'
These ratios will all have the same value: k which is known as the constant of proportionality.
a = k a'
b = k b'
c = k c'
It is this property of similar triangles which enables trigonometry to work.
The ratios between the lengths of the sides of similar triangles always remain the same.
The angle of interest: B is the one between the horizontal and the radius.
The vertical side: b is opposite: o the angle, the horizontal side: a is adjacent: a to the angle and the radius: side c is the hypotenuse: h.
The trigonometric ratios sine, cosine and tangent are the ratios between the sides of the triangle:
sine (B) = o / h
cosine (B) = a / h
tangent (B) = o / h
Also notice: sin B / tan B = (o/h)/(a/h) = o / a = tan B
Jörg Straube
M.Sc. in Computer Science, ETH Zurich (Graduated 1987)
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Author has 6.2K answers and 1.7M answer views
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3y
Draw a right-angled triangle.
There we name the three sides as follows:
- the adjacent = side adjacent to angle x
- the opposite = side opposite to angle x
- the hypotenuse
Now, we have the following definitions
sin(x) = opposite/hypotenuse
cos(x)= adjacent/hypotenuse
tan(x) = opposite/adjacent
In this triangle, we have x = 53.1°
sin = 12/15 = 0.8
cos = 9/15 = 0.6
tan = 12/9 = 1.33333
Draw a right-angled triangle.
There we name the three sides as follows:
- the adjacent = side adjacent to angle x
- the opposite = side opposite to angle x
- the hypotenuse
Now, we have the following definitions
sin(x) = opposite/hypotenuse
cos(x)= adjacent/hypotenuse
tan(x) = opposite/adjacent
In this triangle, we have x = 53.1°
sin = 12/15 = 0.8
cos = 9/15 = 0.6
tan = 12/9 = 1.33333
Related questions
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What actually are sine, cosine and tangent?
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Current Location>Math Formulas>Geometry> Area of a Square
Area of a Square
A square is a regular polygon with four sides. It has four right angles and parallel sides. To calculate the area of a square, multiply the base by itself, which can be expressed as side × side. If a square has a base of length 8 inches its area will be 8 × 8= 64 square inches
Area of a square is given by:
A = a 2
where a =length of side
Perimeter of a square = 4a
Diagonal of a square = (a)(sqrt(2)) or 1.414 (a)
Example 1:Find the area of a square of side length 15 m
Solution:
Area of a square= a 2 = 15 2 = 225 m 2
Example 2:Calculate the area of square, where the square has 35cm side’s length.
Solution:
Area of square is defined by a × a.
Area = 35 × 35
Area = 1225cm
Example 3:What is the area of a square field, if its perimeter is 32 yd?
Solution:
The perimeter of the square field = 32 yd and since the perimeter of a square is given by P = 4s, where s is the length of the side. We can easily determine the length by isolating s from the formula above:
s = P/4 = 32 / 4 = 8 yd
The area of the square field = s × s
Substitute the value of s, we have:
Area = 8 × 8 = 64 yd 2
The area of the square field is therefore 64 yd 2.
Example 4:The side of a square park is 200 m. What will be the cost of grassing it @ $0.5 per sq m?Solution:
What we need to do, is to find the area of the park then multiply the area by the cost per m 2.
Area of the square park = side × side
A = s²
Substitute values and simplify.
A = 200 × 200
A = 40,000 m 2
Area of grassing = area of park = 40,000 sq m.
Cost of grassing = area of grassing × rate per sq meter.
Substitute values we will get:
Cost = 40,000 x 0.5 = $20,000
Therefore, Cost of grassing is $20,000.
Example 5:A square lawn is surrounded by a path 2 m wide around it. If the area of the path is 160 sq m, find the area of the lawn.
Solution:
Given: A square lawn is surrounded by 2 m wide path; area of path is 160 sq m.
To find: Area of lawn.
(Hint: Lawn is surrounded by the path i.e the path is at external edge of lawn. to find area of lawn subtract area of paths from total area)
Let side of the lawn be y, then we have:
External side including path = side of lawn + width of path on both sides.
= y + (2 + 2)
= y + 4
Total area including path = (y + 4) × (y + 4).
= y² + 8y + 16 (i).
And the area of lawn = (side)² = y × y = y² (ii).
Since the area of the path is given (160 m 2), we have:
Area of path = Total area including path - area of lawn.
A = (i) - (ii).
Substitute the given values we will the following equation and by isolating y, we can determine the length of side of the lawn:
160 = (y² + 8y + 16) - y²
160 = y² + 8y + 16 - y²
160 = y² - y² + 8y + 16
160 = 8y + 16
160 - 16 = 8y
144 = 8y
18 = y
Side of lawn = 18 m
Area of the lawn = side × side
A = s²
A = 18 × 18
A = 324 m 2
Hence the area of lawn = 324 m 2.
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Numbers up to 3 Digit
3-digit numbers begin with 100 and end on 999. These numbers consist of 3 digits in which the first digit should be 1 or greater than 1 and the remaining 3 digits can be any number from 0 to 9. Learning 3-digit numbers is the building block for higher-digit numbers. Let us explore more about the importance, formation, and place value of numbers up to 3 digit.
1.What are 3-Digit Numbers?
2.Place Value of 3-Digit Numbers
3.Expanded Form of 3-Digit Numbers
4.Common Mistakes of Numbers up to 3 Digits
5.Operations of Numbers up to 3 Digits
6.FAQs on Numbers up to 3 Digits
What are 3-Digit Numbers?
3-digit numbers are those numbers that consist of only 3 digits. They start from 100 and go on till 999. For example, 673, 104, 985 are 3-digit numbers. It is to be noted that the first digit of a three-digit number cannot be zero because in that case, it becomes a 2-digit number. For example, 045 becomes 45.
Place Value of 3-Digit Numbers
Every three-digit number’s value can be found by checking the place value of each digit. Let us consider the number 243. The first digit at the rightmost position is said to be at units place, so it will be multiplied by 1. Hence, the product is 3 × 1 = 3. Then the second number is 4, and because it is at tens place, it is multiplied by 10. The value, therefore, is 4 × 10 = 40. The third number 2 is at the hundreds place. So 2 is multiplied by 100 and its value is 2 × 100 = 200. Therefore the number is 200 + 40 + 3 = 243.
Decomposing a 3-digit number: In a three-digit number, there are three place values used – hundreds, tens, and units. Let us take one example to understand it better. Here, 465 is a three-digit number and it is decomposed in the form of a sum of three numbers. As 5 is on the units place, 60 is on the tens place and 400 is on the hundreds place.
Significance of Zero in 3-digit numbers: The number zero does not make any contribution to a 3-digit number if it is placed in a position where there are no other non-zero numbers to its left. So how is 303 different from 033 or even 003? In 033, the values are (0 × 100) + (3 × 10) + (3 × 1) = 0 + 30 + 3= 33 which means that the number actually becomes a 2-digit number, i.e., 33, or in the case of 003, it becomes a single-digit number, i.e., 3. In these two examples, the zero does not contribute any value to the number, so the numbers can be expressed as 33 or 3 as well.
Expanded Form of 3-Digit Numbers
The expanded form of a 3-digit number can be expressed and written in three different ways. Consider a three-digit number 457. The number 457 can be written in one form as 457 = (4 × hundreds) + (5 × tens) + (7 × ones). In the second way, the number 457 can be written as 457 = (4 × 100) + (5 × 10) + (7 × 1). And finally the number 457 can be expanded in the form as 457 = 400 + 50 + 7. All the three ways of writing numbers in the expanded form are correct. Writing a 3-digit number in the expanded form helps to know the constituents of the number.
Basically splitting or expanding a 3-digit number helps us to understand more about the 3-digit number. By splitting we know the number of hundreds, tens, and units available in the 3-digit number.
Important Notes on 3-digit Numbers
100 is the smallest 3-digit number and 999 is the greatest 3-digit number.
A 3-digit number cannot start with 0.
10 tens make 1 hundred which is the smallest 3-digit number and 10 hundred make a thousand which is the smallest 4-digit number.
A 3-digit number can also have two zeros but the two zeros should be in the tens place and the units place, for example, 100, 200, 300, 400. It is to be noted that the zeros cannot be in the hundreds place because in that case it becomes a 2-digit number. For example, 067 becomes 67.
Common Mistakes of Numbers up to 3-Digits
Some of the common mistakes are observed while writing or reading a 3-digit number. These mistake in reading and interpreting a 3-digit number is often understood as some other number. In the process of reading, writing, and interpreting a 3-digit number, the place value of the digits should be rightly interpreted. Here we have listed below the three common mistakes often committed by children in writing three-digit numbers.
Misconception 1: Children make mistakes identifying numbers when there is a zero in the units place or tens place. Example: When asked to read 130 and 103, students may get confused. It helps them to model the numbers through Base-10 blocks. That way they can see the ten’s and one’s place value explicitly.
Misconception 2: When asked to write “one hundred twenty-three," students often write 100 first and then attach 23 to it thus ending up with the number “10023” Fact:This misconception arises due to a superficial understanding of place values. Using the base-10 blocks or abacus show children that a digit has different values based on its position.
Misconception 3: Sometimes when asked to form the smallest 3-digit number given three digits that include zero, children place the zero in the left-most position. Fact:This is incorrect. Zero cannot be in the hundreds place if we are creating a 3-digit number. For example: the smallest 3-digit number using all digits of 5, 0, and 7 is 507 and not 057
Operations on Numbers up to 3-Digits
The four arithmetic operations of addition, subtraction, multiplication, and division can be conveniently performed across 3-digit numbers. In the process of performing these arithmetic operations, the place value of the corresponding number should be rightly matched. An error in matching the place value could result in wrong answers. Here we shall look at a simple activity using 3 digit numbers, to help us understand the changing pattern in each of the digits of the hundredth place, ten's place, and unit's place. This activity shall help in a better understanding of the learning needed for the 3 digit numbers.
Get students to skip count by 10 and 100 to build fluency with 3-digit numbers. First, start at 100. Then start from any random 3-digit number like 136.
Help children spot the pattern that when skip counting by 10, the digit in the ones place value does not change. Similarly, when skip counting by 100, the digits in the ones place and tens place does not change.
Use a 100-square grid to build fluency. Let students spot the pattern that moving one row up or down is the same as skip counting by 10. Moving columns (left or right) increases or decreases numbers by 1.
Often children are given three digits and asked to find the largest and smallest number three-digit number using all digits. The trick here is to arrange all digits in descending order to find the largest number.
To find the smallest number, arrange all digits in ascending order. But keep in mind that if zero is one of the digits, it cannot be placed to the left. E.g. Using the digits 7, 3, and 6, the largest number is 763 (digits in descending order) and the smallest number is 367 (digits in ascending order). Using the digits 4, 0, and 8, the largest number will be 840 but the smallest 3-digit number is 408 and not 048.
Smallest 3-Digit Number
The smallest 3-digit number is 100 because its predecessor is 99 which is a two-digit number. 3-digit numbers start from 100 and end on 999.
Greatest 3-Digit Number
The greatest 3-digit number is 999 because its successor is 1000 which is a four-digit number. 3-digit numbers start from 100 and end on 999.
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Examples on 3-Digit Numbers
Example 1: How many 3-digit numbers are there?
Solution:
There are 900 three-digit numbers in all. This can be calculated using the following method.
Step 1: Write the largest and the smallest 3-digit numbers. We know that the largest 3-digit number is 999. The smallest 3-digit number is 100.
Step 2: Find the difference between them. Their difference is, 999 - 100 = 899
Step 3: Add 1 to the difference. This means 899 + 1 = 900. Therefore, there are 900 three-digit numbers in all.
Example 2: Solve the puzzle: Add the smallest 2-digit number to the smallest 1-digit number. Subtract the sum from one less than the greatest 3-digit number.
Solution:
The smallest 2-digit number = 10. The smallest 1-digit number = 1. The sum of these two numbers is 10 + 1 = 11. One less than the greatest 3-digit number is 998. On subtracting 11 from 998, we get. 998 - 11 = 987.
3. Example 3: Find the greatest 3-digit number which is a perfect square.
Solution: The greatest 3-digit number which is a perfect square is 961 because 31 2 = 961.
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Practice Questions on Numbers up to 3-Digits
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FAQs on Numbers up to 3-digits
How Many 3-Digit Numbers are there?
There are a total of 900 three-digit numbers. These include the smallest 3 digit number - 100 to the largest 3 digit number - 999. The numbers beyond these 3-digit numbers are the 4-digit numbers, and the numbers less than the 3-digit numbers are 2-digit numbers.
Which is the Largest 3 Digit Number?
The largest 3-digit number is 999. Adding 1 more to it will make it a 4-digit number, that is, 1000.
What is the Sum of the Three Largest 3 Digit Numbers?
The three largest 3-digit numbers are 997, 998, 999. Their sum is 2994 as 997 +998 +999 = 2994.
What is the Smallest 3-Digit Number?
The number 100 is the smallest 3-digit number. Subtracting 1 from it makes it a 2-digit number. There are a total of 900 three-digit numbers, of which the number 100 is the smallest 3-digit number.
How Many Even 3-Digit Numbers are there?
There are a total of 900 3-digit numbers. Of these half of them are even numbers and the remaining half are odd numbers. Hence there are 900/2 = 450 even 3-digit numbers.
Can a 3-Digit Number have Two Zeros?
A 3-digit number can have two zeros. The two zeros should be in the tens place and the units place. Some of the examples of 3-digit numbers with two zeros are 100, 200, 300, and 400. It should be noted that the hundreds place in a 3-digit number cannot have the number 0 because that will make it a 2-digit number. For example, 098 becomes 98.
Which is the Smallest 3 Digit Number Divisible by 4?
The smallest 3-digit number is 100 and we know that it is divisible by 4 because 100/4 = 25. Therefore, we can say that 100 is the smallest 3-digit number which is divisible by 4.
Which 3 Digit Number has the Most Factors?
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What is h=1/2gt^2?
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1y
The equation h=1 2 g t 2 h=1 2 g t 2 describes the relationship between the height h h of an object in free fall, the acceleration due to gravity g g, and the time t t that the object has been falling.
Variables:
h h: Height (or distance fallen) in meters (m)
g g: Acceleration due to gravity, approximately 9.81 m/s 2 9.81 m/s 2 on the surface of the Earth
t t: Time in seconds (s) that the object has been falling
Explanation:
Free Fall: This equation applies to objects falling freely under the influence of gravity, with no air resistance.
Quadratic Relationship: The height fallen increases with the square of the time,
Continue Reading
The equation h=1 2 g t 2 h=1 2 g t 2 describes the relationship between the height h h of an object in free fall, the acceleration due to gravity g g, and the time t t that the object has been falling.
Variables:
h h: Height (or distance fallen) in meters (m)
g g: Acceleration due to gravity, approximately 9.81 m/s 2 9.81 m/s 2 on the surface of the Earth
t t: Time in seconds (s) that the object has been falling
Explanation:
Free Fall: This equation applies to objects falling freely under the influence of gravity, with no air resistance.
Quadratic Relationship: The height fallen increases with the square of the time, meaning if you double the time, the height fallen will increase by a factor of four.
Initial Velocity: The equation assumes that the initial velocity of the object is zero (i.e., it starts from rest).
Example:
If an object falls for 3 seconds, the height fallen can be calculated as follows:
h=1 2⋅9.81 m/s 2⋅(3 s)2=1 2⋅9.81⋅9=44.145 m h=1 2⋅9.81 m/s 2⋅(3 s)2=1 2⋅9.81⋅9=44.145 m
Thus, the object would fall approximately 44.145 meters in 3 seconds.
Upvote ·
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More answers below
Why do we use 1/2gt2 to find height?
How do you evaluate the free fall formula (h=1/2gt^2)?
How do you prove H=1/8gt^2?
If h=ut +1/2gt^2 then what is the dimension of g?
What is h in the equation: h=ut+1/2gt^2? Is h the distance between the object and the topmost position or the distance between the object and the ground?
C Isaac
Secondary Diploma from Po Leung Kuk Centenary Li Shiu Chung Memorial College (Graduated 2023)
·5y
It is one of the formulas in physics.
h=displacement (height)
g=9.8 ms^-2 (Cor.to 1 d.p)=constant of the acceleration of gravitational force=the change in speed of the object through the motion of falling
t=time (in seconds)
It calculates the motion of an object falling from a higher point to a lower point through a certain distance(or height) h within a certain period of time t, plus the influence of gravity g exerted from the earth.
However, the following formula is used more commonly:
s=ut+1/2 (at^2)
where s=displacement, a=9.8, t= time (sec), u=initial velocity(speed) of the object(ms^-1 or m/s).
Continue Reading
It is one of the formulas in physics.
h=displacement (height)
g=9.8 ms^-2 (Cor.to 1 d.p)=constant of the acceleration of gravitational force=the change in speed of the object through the motion of falling
t=time (in seconds)
It calculates the motion of an object falling from a higher point to a lower point through a certain distance(or height) h within a certain period of time t, plus the influence of gravity g exerted from the earth.
However, the following formula is used more commonly:
s=ut+1/2 (at^2)
where s=displacement, a=9.8, t= time (sec), u=initial velocity(speed) of the object(ms^-1 or m/s).
Moreover, the formula was actually founded by using calculus.
In calculus,
therefore,
then
q.e.d.
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Altaha Ansari
B.tech(ECE), Ambivert, Student, Sometimes 0 Sometimes 1
·6y
This is one of the euations of motion i.e. S=ut + 1/2at²
Here S is displacement, U is initial Velocity, a is acceleration, t is time.
When a body falls under gravity when it's initial velocity is zero and acceleration is due to gravity and g remains fairly constant when height h is not large, the equation reduces to the form h= 1/2gt².
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99 43
9 6
Calvin Campbell
B.Sc. in Computer Science&Mathematics, University of the West Indies (Graduated 1984) · Author has 3.9K answers and 4.1M answer views
·Updated 5y
This is a formula from elementary physics, where the (h) is the distance a projectile has fallen after time (t), assuming a starting velocity of zero. The constant (g) is acceleration due to gravity and has a value of 9.8 meters per second squared (9.8m/s^2).
Examples: If you drop a ball from a 500 meters cliff, in 10 seconds it would fall a distance, h = 1/29.8m/s^2(10s)^2 = 490 meters .
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How big of a generator do I need to run a well pump that is 1/2 horsepower and 6 amps?
How does s=ut+1/2 at^2?
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What is 3+33-3+3?
How many amps does a 1/2 horsepower sump pump draw?
Jacob Valdehueza
theoretical physics in Physics&Mechanical Engineering, Harvard University (Graduated 2017)
·7y
damn it! it’s just a formula to calculate the height of an object.
it’s based on vertical motions of an object which derived accordingly by any questions.
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99 65
Gary Ward
MaEd in Education&Mathematics, Austin Peay State University (Graduated 1997) · Author has 4.9K answers and 7.6M answer views
·4y
Related
Why do we use 1/2gt2 to find height?
Why do we use ½gt² to find height?
A body in freefall, in a vacuum , in a close to uniform gravitational field will have a constant acceleration, the rate of change in velocity. If you drop a bowling ball off the top of a ladder, the air drag can be ignored for practical purposes and the gravitational field is uniform enough that we can ignore the slight change. In truth, all objects in our atmosphere will experience a decrease in the rate of increase of their velocity as they approach their terminal velocity where the force of air drag equals the force of gravity acting upon the object.
½gt²=s,
Continue Reading
Why do we use ½gt² to find height?
A body in freefall, in a vacuum , in a close to uniform gravitational field will have a constant acceleration, the rate of change in velocity. If you drop a bowling ball off the top of a ladder, the air drag can be ignored for practical purposes and the gravitational field is uniform enough that we can ignore the slight change. In truth, all objects in our atmosphere will experience a decrease in the rate of increase of their velocity as they approach their terminal velocity where the force of air drag equals the force of gravity acting upon the object.
½gt²=s, where ‘s’ is the height will give a negative value since the distance is from zero height to a lower value. The value of ‘g’ should be -9.8 m/s² in the mks system and -32 ft/s² in the conventional system. It simply means that for every passing second the downward velocity increases by that amount.
Distance, s, is s = v · t, where v is the velocity. For a velocity changing at a uniform rate, we need the average velocity, (v1 - v0) / 2, where v1 is the final velocity and v0 is the starting velocity. To simplify, if the starting velocity is zero, the average velocity is (v1)/2. The final velocity, v1 = a · t, where a is the acceleration. Substitute a · t for v1 and get the average velocity is (a · t)/2 = ½(a · t).
Substitute ½(a · t) into the equation for v. s = v · t becomes s = ½(a · t) · t = ½at². g=a, so s = ½gt²
Usually you will see this used with a constant ,c, s = c +½gt², where a question might read, “A ball rolls off a 100 meter cliff. How long will it take to hit the ground?”
s = c +½gt² becomes s = 100 m+ ½(-9.8 m/s²t²) The height is zero at the bottom of the cliff, so s = 0. 0 = 100 m+ ½(-9.8 m/s²t²)
t² = 200 m / 9.8 m/s²= 20.4 s² means t = 4.52 s
“What was its height after 2 seconds?”
s = 100 m+ ½(-9.8 m/s²2²) becomes s = 100 m+ -19.6 m becomes s = 80.4 m
What if the ball was thrown upwards or downwards with some initial velocity?
Then the whole equation becomes s = c+ vt -4.9 t² in its stripped down form where vt is positive if thrown upward and negative if thrown downward,
You can see it is a nice parabola.
“How fast was it moving when it hit the ground?”
v = a · t + v0 becomes v = -9.8 m/s² · 4.52 s + 0 = -44.3 m/s where the negative sign indicates a downward direction.
That’s as simple as I can present it without using calculus.
Upvote ·
9 6
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9 1
Charlie Thiery
B.S. in Electrical Engineering, University of Illinois at Urbana-Champaign (Graduated 2022)
·7y
It is the formula for how far you have fallen for a specific time in freefall with a starting velocity of 0 m/s.
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Allen Ries
Math Major University of Alberta · Author has 25.1K answers and 9.7M answer views
·1y
Related
What is h in the equation: h=ut+1/2gt^2? Is h the distance between the object and the topmost position or the distance between the object and the ground?
What is h in the equation: h=ut+1/2gt²? Is h the distance between the object and the topmost position or the distance between the object and the ground?
It is NOT the the distance between the object and the topmost position.
It is NOT the distance between the object and the ground.
h measures the vertical distance of the object at time ‘t’ COMPARED to the initial position of the object when the time ‘t’ is zero.
It is important that the vertical position ‘h’ be zero when the time ‘t’ equals zero.
Upvote ·
9 2
Abdur Rahaman Srabon
BSc. in Electrical and Electronics Engineering, Islamic University of Technology
·6y
Here h=Height traveled by a projectile after time ‘t’ ;
g= Acceleration due to gravity ;
t=Time spent to gain height ‘h’ .
This equation can be used to get any unknown h,g or t if another 2 variable is given.
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Maneet Jha
Studied Physics, Chemistry, and Mathematics (science grouping) at FIITJEE - EAST DELHI (Graduated 2019)
·6y
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What is h in the equation: h=ut+1/2gt^2? Is h the distance between the object and the topmost position or the distance between the object and the ground?
As u have mentioned 'g' there which means the object is in freefall under the effect of gravity which means object's acceleration is equal to g which is 9.8 m/s^2
Now,Coming to this kinematics equation
Here,it says if object is with a initial velocity 'u'
And travelles with a acceleration for time t then distance covered by it will be h
And coming to your answer as object falls from a point take it is topmost point and suppose falls on anything or on ground then distance between both point will be 'h'
So,yes you are correct It is distance between topmost point and point where object stopped.
I hope
Continue Reading
As u have mentioned 'g' there which means the object is in freefall under the effect of gravity which means object's acceleration is equal to g which is 9.8 m/s^2
Now,Coming to this kinematics equation
Here,it says if object is with a initial velocity 'u'
And travelles with a acceleration for time t then distance covered by it will be h
And coming to your answer as object falls from a point take it is topmost point and suppose falls on anything or on ground then distance between both point will be 'h'
So,yes you are correct It is distance between topmost point and point where object stopped.
I hope it is clear,
Feel free to comment
Thanks
Upvote ·
9 7
Colin Watters
BSC in Electrical and Electronics Engineering&Ex Glider Pilot, University of Southampton (Graduated 1981) · Author has 4.7K answers and 2.9M answer views
·1y
Related
What is h in the equation: h=ut+1/2gt^2? Is h the distance between the object and the topmost position or the distance between the object and the ground?
It's the distance between its starting position at t=0 and some time t.
So it could be either the distance from the top or from the ground depending on where it started from at t=0.
Upvote ·
9 2
Richard Chiang
Former Retired (1973–2002) · Author has 714 answers and 473.3K answer views
·4y
Related
Why do we use 1/2gt2 to find height?
Height is distance travelled by a falling object or a projectile.
Velocity = Distance/Time or Distance = Velocity x Time
In the case of a falling object, the initial velocity is zero. Hence, average velocity is
= (initial velocity + final velocity)/2 = (0 + v)/2 = v/2.
Distance, d = v/2 x t ……………………………………………………………………….A
Now, acceleration is the rate of change of velocity of an object with respect to time.
Hence, acceleration, a = Delta v/Delta t = v/t. => v = at ……..……B
Combining equations A and B, d = 1/2at^2
Acceleration due to gravity (for a falling object) is g = approximately 9.8m/sec^2.
Hence, d
Continue Reading
Height is distance travelled by a falling object or a projectile.
Velocity = Distance/Time or Distance = Velocity x Time
In the case of a falling object, the initial velocity is zero. Hence, average velocity is
= (initial velocity + final velocity)/2 = (0 + v)/2 = v/2.
Distance, d = v/2 x t ……………………………………………………………………….A
Now, acceleration is the rate of change of velocity of an object with respect to time.
Hence, acceleration, a = Delta v/Delta t = v/t. => v = at ……..……B
Combining equations A and B, d = 1/2at^2
Acceleration due to gravity (for a falling object) is g = approximately 9.8m/sec^2.
Hence, d = 1/2gt^2.
Upvote ·
9 2
Fakhri As'Ad
Studied Physics (Graduated 1979) · Author has 3.7K answers and 814.7K answer views
·1y
Related
What is h in the equation: h=ut+1/2gt^2? Is h the distance between the object and the topmost position or the distance between the object and the ground?
h is the distance traveled by the object in t time between the reference frame and the object, the reference frame could be ground or the starting point.
Upvote ·
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13765 | https://physics.stackexchange.com/questions/556745/how-would-i-explain-to-a-child-that-zero-to-the-zeroth-power-equals-one | mathematics - How would I explain to a child that zero to the zeroth power equals one? - Physics Stack Exchange
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How would I explain to a child that zero to the zeroth power equals one? [closed]
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I’m trying to figure out why raising zero to the zeroth power equals one. What kind of a scenario would occur in a laboratory experiment where something with a quantity of zero would be raised to the power of zero and you end up with one? How do I explain how something is created out of nothing? What is happening?
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edited Jun 3, 2020 at 3:51
Qmechanic♦
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asked Jun 3, 2020 at 0:56
Thom Blair IIIThom Blair III
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5 perhaps better asked on Mathematics? I don't think you can relate the 0^0=1 issue to a lab experiment, it is a feature of the math.user9413641 –user9413641 2020-06-03 01:02:26 +00:00 Commented Jun 3, 2020 at 1:02
10 Maybe you shouldn't explain that zero to the power of zero is one because it isn't necessary: en.wikipedia.org/wiki/Zero_to_the_power_of_zerouser87745 –user87745 2020-06-03 01:08:48 +00:00 Commented Jun 3, 2020 at 1:08
@DvijD.C. Thank you. I didn’t realize it was undefined. I appreciate the help.Thom Blair III –Thom Blair III 2020-06-03 02:40:54 +00:00 Commented Jun 3, 2020 at 2:40
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The thing is, zero to the zeroth power is not always one. It is technically undefined. We can try to get around that by asking the question; What happens if we take the limit of a function that approaches zero to the zeroth power? It should be easy to find two functions that both approach zero at some point (lets call them f f and g g). Thus, f f to the power of g g approaches zero to the zero. However, using L'hospital's rule, we can see that it is possible to find f f and g g so that f g f g approaches a couple of different numbers (all real numbers, if you include the complex numbers)
That will probably be too technical for your child, so to explain it to them, I would recommend asking them what zero to the power of <<any number besides zero>> is, and showing them that it is always zero. Then, ask them what <<any number besides zero>> to the zeroth power is, and show them that it is always one. Then explain that because these two properties conflict, there is no single right answer. To make it easier to understand, it is standard practice to say that it equals one, even though it is technically incorrect.
Edit: I looked through my explanation of the math behind the weirdness, and found a number of errors. My math was wrong, and the example I gave didn't actually fit the requirements. I removed those parts of my answer. I believe that the simplified explanation for a child is still accurate, though.
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edited Jun 3, 2020 at 2:46
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answered Jun 3, 2020 at 1:52
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Yes, the indeterminate form cannot be explained to a small child. Perhaps, it should be hinted at.Ishika_96_sparkle –Ishika_96_sparkle 2020-06-03 01:54:43 +00:00 Commented Jun 3, 2020 at 1:54
Thank you to whoever suggested the numberphile video; it was very helpful, but the comment was removed, and I don't know who posted it. Here is the link for anyone else who wants to see it: youtube.com/watch?v=BRRolKTlF6QHyrum Taylor –Hyrum Taylor 2020-06-03 02:41:40 +00:00 Commented Jun 3, 2020 at 2:41
x 0→1 x 0→1 when x→0 x→0. This is perfectly well defined.Charles Francis –Charles Francis 2020-06-03 08:30:47 +00:00 Commented Jun 3, 2020 at 8:30
0 y 0 y is not defined at y≤0 y≤0. This has no bearing.Charles Francis –Charles Francis 2020-06-03 08:42:47 +00:00 Commented Jun 3, 2020 at 8:42
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If we consider this operation (a≠0 a≠0)
a n a n=a n−n=a 0=1 a n a n=a n−n=a 0=1
and taking a special case for illustration with a=e a=e i.e. the exponential, we could write
e 0=1 e 0=1
and then taking natural log on both sides
0 ln e=ln 1=0 0 lne=ln1=0
If we now replace a by 0 0, then
0 n 0 n=0 n−n=0 0=1 0 n 0 n=0 n−n=0 0=1
It seems to suggest that zero is quite fundamental; as it gives something out of nothing when not multiplied with itself!
However, the 0 0 0 0 is indeterminate in calculus.
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edited Jun 3, 2020 at 2:01
answered Jun 3, 2020 at 1:53
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you cannot divide by zero!user65081 –user65081 2020-06-03 01:58:05 +00:00 Commented Jun 3, 2020 at 1:58
Yes, that is why in calculus the form 0 0 0 0 is indeterminate. If one cannot divide zero by zero then 0 0 0 0 cannot happen.Ishika_96_sparkle –Ishika_96_sparkle 2020-06-03 01:59:53 +00:00 Commented Jun 3, 2020 at 1:59
so your last equation is undefined, not equal to one user65081 –user65081 2020-06-03 02:00:59 +00:00 Commented Jun 3, 2020 at 2:00
Yes, that is why i wrote -it seems to suggest.Ishika_96_sparkle –Ishika_96_sparkle 2020-06-03 02:02:04 +00:00 Commented Jun 3, 2020 at 2:02
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In order to explain it, you need to have them to have some clear idea of what "raising to a power" means, because "explanation" itself fundamentally means you are trying to describe something as the result of some more basic principle.
Hence, what we're really talking about is how to best define exponentiation. And in that regard, we should set it up by saying that there are at least two possible avenues to defining exponentiation at this level. Most likely, your kid will think of exponentiation as being something like this:
a n:=a⋅a⋅⋯⋅an appearances of "a"a n:=a⋅a⋅⋯⋅a⏟n appearances of "a"
and if you take this literally and try to put n=0 n=0, you get
a 0=see, there aren't any copies of "a" here"!a 0=⏟see, there aren't any copies of "a" here"!
which is nonsense, of course. Hence this definition fails to define a value for a 0 a 0 - whether a a is zero or not.
The problem is, we are defining exponentiation by the idea that it simply represents the "number of 'copies'" of a a that appear in some thing we write down on paper that in just some instances happens to correspond to a valid mathematical expression. But clearly, that's problematic, and your kid (and you and perhaps many others) have figured that out. So, we need to do something else, which motivates us to look out for a different approach.
And that's this. Instead of thinking of a n a n is being defined by putting "copies" or "appearances" of a a into some expression, think of n n as saying "how many times do we multiply some fundamental base value by a a"? That is,
a n:=K⋅a⋅a⋅a⋯⋅an performances of the operation "⋅a"a n:=K⋅a⋅a⋅a⋯⋅a⏟n performances of the operation "⋅a"
(note the text! We're talking about performing an operation as what is counted.)
If n=0 n=0, then this should just give back the fundamental base value, K K, instead of giving us a nonsensical expression, because we will have not performed the multiplication operation at all. That is, a 0=K a 0=K. Hence the value of a 0 a 0 depends wholly on the choice of initial constant K K, and not the structure of the operation.
But of course, we still need to make such a choice, as the definition is thus incomplete! Hence, what value we should set for K K? And that comes to this. If they are into the idea of multiplication as being interpretable as a scaling or "zooming" process, then you can say that the motivation for the operation a n a n is to tell us what "zoom" factor performs n n successive "zooms" of magnification a a. For n=0 n=0, it should mean we do no zooms, which is a zoom factor of 1 1. Hence we should define K:=1 K:=1, and the fully available definition for exponentiation looks like:
a n:=1⋅a⋅a⋅a⋯⋅an performances of the operation "⋅a"a n:=1⋅a⋅a⋅a⋯⋅a⏟n performances of the operation "⋅a"
where we purposefully do not just elide the "1 1" - something too many miss perhaps in part thanks to being taught that so-called "simplifying" is a "virtue". (If anything, I say to ditch "simplifying" in favor of clarifying.)
And it then follows from this that
0 0=1 0 0=1
and 0 0 is not more special than any other number that can go into the exponentiation. If you want to think of it as "zooming", a zoom factor of 0 0 "crushes everything into a singularity" (as zero annihilates everything it's multiplied with). Yet, no matter how powerful or all-consuming something is, if we never do it in the first place, it doesn't matter. Hence, so long as we do n=0 n=0 of those all-crushing zooms, things should still stay the way they are, i.e. zoom factor 1 1.
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edited Jun 3, 2020 at 2:09
answered Jun 3, 2020 at 2:04
The_SympathizerThe_Sympathizer
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If you multiply a number, N N, by zero no times, then you do nothing to the number, which is the same as multiplying by 1 1.
0 0 N=N=1 N 0 0 N=N=1 N
Hence 0 0=1 0 0=1 (contrary to other answers, this is well defined).
You may also consider that x 0→1 x 0→1 as x→0 x→0, but it is simpler for a child to recognise that power of 0 0 means do no multiplications.
0 0 0 0 may not occur in the lab, but it will occur in the formulae used to describe physics. It is still defined and essentially refers to the case where nothing is done.
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edited Jul 18, 2020 at 7:12
answered Jun 3, 2020 at 8:19
Charles FrancisCharles Francis
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1 I don't get how this works: if you do nothing to a number, then how does that imply 1? It's just interlinking different facts of mathematics without context.PNS –PNS 2020-07-17 14:28:43 +00:00 Commented Jul 17, 2020 at 14:28
@PNS, What is there to get? Multiplying by 1 does nothing to a number. This is a central idea in mathematics and is what makes the number 1 particularly important. 1 is known as the identity under multiplication.Charles Francis –Charles Francis 2020-07-17 16:31:36 +00:00 Commented Jul 17, 2020 at 16:31
Yes, but if you take a number and do nothing, you are multiplying by 1, so from what point does that mean that mean you get only 1 and not the number itself? (From your argument in the answer, you should get 0, as 1 multiplied by 0, is just 0.)PNS –PNS 2020-07-18 03:20:03 +00:00 Commented Jul 18, 2020 at 3:20
@PNS, the argument already says you get the number itself. It does not say you get 1, and it does not multiply 1 by 0. It was already clear, but I have written the multiplication explicitly. I hope that helps.Charles Francis –Charles Francis 2020-07-18 07:17:36 +00:00 Commented Jul 18, 2020 at 7:17
Yes, but now you assume that 0 0 0 0 is working on something N N and you get 1 N 1 N. What if you remove N N now? You don't get 1 1. The 0 0 0 0 is not acting as an 'operation' now, so you don't even get 1 1 anymore.PNS –PNS 2020-07-18 10:41:12 +00:00 Commented Jul 18, 2020 at 10:41
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How to Multiply Multiples of 10: A Step-by-Step Guide for 3rd Graders
Try Our AI Tutor to Learn Multiply Multiples of 10!
Want to make multiplying with 10 super easy? The AI Tutor is here to help! When you multiply a number by 10, you’re just adding a zero at the end. It’s a simple math trick that helps with larger problems later on.
With the AI Tutor, you’ll get step-by-step help, fun practice questions, and instant feedback. It’s designed for 3rd grade students who want to get stronger in math while learning at their own pace. You’ll understand how multiplying by 10 works, and why it’s useful!
Here’s a sample question you can try right now:
Multiply:
4 x 10 = ____
A. 40
B. 10
C. 100
D. 4
Practice Now!
Want more help from the AI tutor?Get Free Guided Practice!
Table of Contents
What Are Multiples of 10?
Why Is Multiplying by Multiples of 10 Important?
Key Concept: Place Value Pattern
Multiplying by Other Multiples of 10
Step-by-Step: How to Multiply Multiples of 10
Let’s Practice With More Examples!
Word Problems Using Multiples of 10
Tips and Tricks for Students to Learn Multiply Multiples of 10
Challenge Questions (For Extra Practice)
Real-Life Applications on Multiplying by multiples of 10
Additional Resources to Learn Multiply Multiples of 10
Final Thoughts
Multiplying numbers is one of the most important skills in 3rd grade math. Once students learn the basics of multiplication, they move on to learn how to multiply multiples of 10. This is a key skill that builds the foundation for understanding place value, mental math, and larger multiplication problems in later grades.
In this guide, we’ll explore what multiples of 10 are, why they matter, how to multiply by them easily, and how this skill applies to real-life situations. With clear explanations, step-by-step instructions, examples, and fun tips, this article is perfect for helping 3rd graders understand how to multiply by multiples of 10.
What Are Multiples of 10?
Let’s start with the basics.
A multiple of 10 is a number that you get when you multiply 10 by another whole number.
Examples of Multiples of 10:
10 × 1 = 10
10 × 2 = 20
10 × 3 = 30
10 × 4 = 40
10 × 5 = 50
10 × 6 = 60
10 × 7 = 70
10 × 8 = 80
10 × 9 = 90
10 × 10 = 100
So, 10, 20, 30, 40, 50, and so on are all multiples of 10 because they can be evenly divided by 10.
Key Clue: All multiples of 10 end in zero (0)!
Why Is Multiplying by Multiples of 10 Important?
Multiplying by multiples of 10 helps students develop strong number sense and understand how our number system works.
Here’s why it’s such an important skill in 3rd grade:
It introduces place value patterns.
It builds mental math strategies.
It makes large multiplication problems easier to solve.
It connects to money, math, measurement, and time.
It prepares students for multi-digit multiplication in 4th grade.
Students who can confidently multiply by 10, 20, 30, and other multiples will be ready for more advanced math problems!
Key Concept: Place Value Pattern
Multiplying by multiples of 10 relies heavily on the idea of place value.
Example 1:
What is 4 × 10?
We know: 4 × 10 = 40
Here’s the trick:
Start with 4, then add a 0 at the end → 40
Why does this work?
Because multiplying by 10 shifts the number one place to the left in our number system, making it 10 times bigger.
Example 2:
What is 7 × 100?
7 × 100 = 700
This time, we added two zeros because 100 has two zeros.
The number of zeros in the multiple of 10 tells us how many place value positions the number shifts.
Multiplying by Other Multiples of 10
It’s not just 10! We can multiply numbers by 20, 30, 40, 50, and even larger multiples of 10 using similar strategies.
Example 3:
What is 6 × 30?
Break it down:
Think of 30 as 3 × 10
So, 6 × 30 = 6 × (3 × 10)
Group it: (6 × 3) × 10 = 18 × 10 = 180
So, 6 × 30 = 180
This method works for any multiple of 10. Just multiply the non-zero numbers first, then add a zero at the end.
More Examples:
8 × 40 = (8 × 4) × 10 = 32 × 10 = 320
5 × 60 = (5 × 6) × 10 = 30 × 10 = 300
9 × 70 = (9 × 7) × 10 = 63 × 10 = 630
Step-by-Step: How to Multiply Multiples of 10
Step 1: Look at the numbers you’re multiplying
Let’s say the problem is 30 × 7
One of those numbers (30) is a multiple of 10.
Step 2: Take off the zero from the multiple of 10
30 has one zero, so we’ll ignore it for now. That leaves us with 3.
Now we multiply: 3 × 7
Step 3: Multiply the digits
3 × 7 = 21
Step 4: Add the zero back
Since 30 had one zero, just add it back to 21 → 210
So, 30 × 7 = 210
Example 4:
Multiply 4 × 50
Step 1: 50 has a zero, so ignore it for now → 5
Step 2: 4 × 5 = 20
Step 3: Add the zero back → 200
Solution: 4 × 50 = 200
This method works like magic!
Let’s Practice With More Examples!
Example 1:
3 × 10 = ?
3 × 10 = 30
Example 2:
7 × 40 = ?
7 × 4 = 28,
add a zero → 280
Example 3:
5 × 70 = ?
5 × 7 = 35,
add a zero → 350
Example 4:
8 × 20 = ?
8 × 2 = 16,
add a zero → 160
Example 5:
6 × 90 = ?
6 × 9 = 54,
add a zero → 540
Word Problems Using Multiples of 10
Problem 1:
Lila has 5 packs of pencils. Each pack has 10 pencils. How many pencils does she have in total?
5 × 10 = 50 pencils
Problem 2:
There are 4 buses. Each bus carries 30 students. How many students are there in total?
4 × 30 = 4 × 3 = 12 → Add a zero → 120 students
Problem 3:
Tom buys 7 notebooks. Each notebook costs $20. How much money did he spend?
7 × 20 = 7 × 2 = 14 → Add a zero → $140
Problem 4:
A box holds 60 crayons. If a class has 6 boxes, how many crayons are there?
6 × 60 = 6 × 6 = 36 → Add a zero → 360 crayons
Practice Problems
Try solving these. (Answers at the bottom!)
Questions:
9 × 10 = ___
6 × 30 = ___
3 × 70 = ___
8 × 90 = ___
2 × 60 = ___
5 × 50 = ___
4 × 40 = ___
7 × 20 = ___
10 × 10 = ___
1 × 80 = ___
Answers:
90
180
210
720
120
250
160
140
100
80
Tips and Tricks for Students to Learn Multiply Multiples of 10
Here are some fun and helpful tips for students to remember:
Tip 1: Look for the Zero – Multiples of 10 always end in zero, so it’s your best friend in these problems.
Tip 2: Break Down the Problem – Think of 60 as 6 × 10, 90 as 9 × 10, and so on. That makes it easier to solve.
Tip 3: Use Arrays or Base-10 Blocks – Drawing groups of tens or using base-10 blocks helps students understand the concept visually.
Tip 4: Practice with Skip Counting – Skip count by 10s to find the answer:
Example: 4 × 10? Count: 10, 20, 30, 40 → Answer: 40
Tip 5: Use a Place Value Chart – Seeing how digits shift in a place value chart when you multiply by 10 or 100 helps kids see the pattern.
Challenge Questions (For Extra Practice)
Challenge 1:
If a farmer has 9 baskets, and each basket holds 70 apples, how many apples does he have in total?
9 × 70 = 9 × 7 = 63 → Add a zero → 630 apples
Challenge 2:
There are 6 shelves in a library. Each shelf holds 90 books. How many books are there in total?
6 × 90 = 6 × 9 = 54 → Add a zero → 540 books
Challenge 3:
If a toy costs $30, how much would 8 toys cost?
8 × 30 = 8 × 3 = 24 → Add a zero → $240
Real-Life Applications on Multiplying by multiples of 10
Multiplying by multiples of 10 isn’t just for the classroom. We use it in everyday life!
Money: You buy 3 toys that cost $30 each → 3 × 30 = $90
Exercise: You do 5 sets of 20 jumping jacks → 5 × 20 = 100 jumping jacks
Transportation: A bus seats 40 people. If there are 6 buses, how many people can ride?
→ 6 × 40 = 240 people
4. Shopping: A box holds 10 cans. If you buy 8 boxes, how many cans do you have?
→ 8 × 10 = 80 cans
Additional Resources to Learn Multiply Multiples of 10
Looking for more practice? Check out these engaging, printable multiplication worksheets and sample problems designed to help 3rd-grade students master multiplication by multiples of 10:
Final Thoughts
Multiplying by multiples of 10 is a simple yet powerful skill that every 3rd grader should master. It makes larger math problems easier and teaches students important place value patterns. Whether you’re using this in class, at home, or out in the world, knowing how to multiply by 10, 20, 30, and more will give you a strong math foundation.
Keep practicing, use visuals like number lines and base-10 blocks, and remember the trick: multiply the numbers, then add the zero!
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13769 | https://math.stackexchange.com/questions/926598/proof-area-of-a-cyclic-quadrilateral | geometry - Proof - Area of a cyclic quadrilateral - Mathematics Stack Exchange
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Proof - Area of a cyclic quadrilateral
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So we have a cyclic quadrilateral, as depicted below:
I have a conjecture that the area of this cyclic quadrilateral equals
(a+b+c−d)(a+b+d−c)(a+c+d−b)(b+c+d−a)−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−−√4(a+b+c−d)(a+b+d−c)(a+c+d−b)(b+c+d−a)4
I want to prove this. I know that the area of triangle ABC equals 1 2 a b sin(B)1 2 a b sin(B) and the area of triangle ACD equals 1 2 c d sin(D)1 2 c d sin(D). Seeing as sin(B)=sin(D)sin(B)=sin(D), I figured that the area of the quadrilateral equals 1 2(a b+c d)sin(B)1 2(a b+c d)sin(B).
But I'm stuck here. I have no clue how to get to the result I want. Can anyone provide an insight (and tell me whether or not my work up until now is valid)?
geometry
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edited Sep 10, 2014 at 17:54
PhaptitudePhaptitude
asked Sep 10, 2014 at 17:48
PhaptitudePhaptitude
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2 Your conjecture is close but wrong. en.wikipedia.org/wiki/Brahmagupta%27s_formulamathlove –mathlove 2014-09-10 17:52:40 +00:00 Commented Sep 10, 2014 at 17:52
@mathlove I edited it. I was iffy about that part, but of course the main problem is the numerator.Phaptitude –Phaptitude 2014-09-10 17:55:09 +00:00 Commented Sep 10, 2014 at 17:55
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The Inscribed Angle Theorem shows that opposite angles of a cyclic quadrilateral are supplementary. Thus,
D=π−B D=π−B
Therefore, cos(D)=−cos(B)cos(D)=−cos(B) and sin(D)=sin(B)sin(D)=sin(B).
The Law of Cosines says
e 2=a 2+b 2−2 a b cos(B)as a side of△A B C=c 2+d 2+2 c d cos(B)as a side of△A D C e 2=a 2+b 2−2 a b cos(B)⏞as a side of△A B C=c 2+d 2+2 c d cos(B)⏞as a side of△A D C
Solving for cos(B)cos(B), we get
cos(B)=1 2 a 2+b 2−c 2−d 2 a b+c d cos(B)=1 2 a 2+b 2−c 2−d 2 a b+c d
Furthermore,
1 4 sin 2(B)=(2 a b+2 c d)2−(a 2+b 2−c 2−d 2)2 16(a b+c d)2=((c+d)2−(a−b)2)((a+b)2−(c−d)2)16(a b+c d)2=((c+d+a−b)(c+d+b−a))((a+b+c−d)(a+b+d−c))16(a b+c d)2=(s−a)(s−b)(s−c)(s−d)(a b+c d)2 1 4 sin 2(B)=(2 a b+2 c d)2−(a 2+b 2−c 2−d 2)2 16(a b+c d)2=((c+d)2−(a−b)2)((a+b)2−(c−d)2)16(a b+c d)2=((c+d+a−b)(c+d+b−a))((a+b+c−d)(a+b+d−c))16(a b+c d)2=(s−a)(s−b)(s−c)(s−d)(a b+c d)2
where s=a+b+c+d 2 s=a+b+c+d 2 is the semi-perimeter.
The area of ⏥A B C D⏥A B C D is the sum of the area of △A B C=1 2 a b sin(B)△A B C=1 2 a b sin(B) and the area of △A D C=1 2 c d sin(D)△A D C=1 2 c d sin(D)
1 2(a b+c d)sin(B)=(s−a)(s−b)(s−c)(s−d)−−−−−−−−−−−−−−−−−−−−−−√1 2(a b+c d)sin(B)=(s−a)(s−b)(s−c)(s−d)
which, as Christian Blatter mentions, is called Brahmagupta's Formula.
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edited Apr 28, 2015 at 15:20
answered Apr 27, 2015 at 12:16
robjohn♦robjohn
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Why this remained without any upvote for so long is a mystery. +1 Paramanand Singh –Paramanand Singh♦ 2017-12-07 08:48:08 +00:00 Commented Dec 7, 2017 at 8:48
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These might be useful:
Heron's formula:
Δ=s(s−a)(s−b)(s−c)−−−−−−−−−−−−−−−−−√Δ=s(s−a)(s−b)(s−c)
Ptolemy's theorem:
A C.B D=a c+b d A C.B D=a c+b d
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answered Sep 10, 2014 at 17:54
RE60KRE60K
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This is Brahmagupta's formula; see here:
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answered Sep 10, 2014 at 18:56
Christian BlatterChristian Blatter
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5What is the expected area of a cyclic quadrilateral inscribed in a unit circle?
2How do I prove that the area of any triangle can be obtained using "Heron’s Formula" ⟶A t=s(s−a)(s−b)(s−c)−−−−−−−−−−−−−−−−−√⟶A t=s(s−a)(s−b)(s−c)?
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12Prove that a quadrilateral, and the quadrilateral formed by the orthocenters of four related triangles, have the same area.
1Cyclic quadrilateral, and isogonal conjugates?
1Area of Inscribed (Cyclic) Quadrilateral
3Area of a cyclic quadrilateral.
1Solving area of quadrilateral from given values
0An obscure problem related to cyclic quadrilaterals
1Quadrilateral A B C D A B C D, with A B=A C A B=A C, C D⊥B C C D⊥B C, S A B C D=54 S A B C D=54 and S△A C D=17 S△A C D=17. What is the area of △A B D△A B D?
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13770 | https://www.youtube.com/watch?v=IyPFw4aJUwI | Why is X^0 = 1 ?? | Exponents Proof | Mathematics | dArtofScience
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Today we will tap on yet another mathematical proof and find out why any number raised to the power 0 = 1??
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Intro [Music] hey guys I'm channi and you're watching the art of science so in today's video we are going to do yet another mathematical proof and don't worry guys it's not a fake proof this time it's a very real one so when I was in school and we were learning exponence one of the things that you learn is any number raised to the power 0 is equal to 1 and this is something that I never understood shouldn't a number rais to the^ 0 be equal to 0 how is it equal to 1 so today I'm going to make a video exactly on that this is the proof of why X to the^ 0 is equal to one I hope that all the Curious Kids wondering the same question stumble across my video and find the answer to your questions so let's get into it so before we get into What are exponents the proof we are going to first understand a little is something about exponents exponents are usually written something like this the number in the bottom is called the base and the smaller one on the top right side is called the power here we have 2 to the^ 3 this basically means that the base which is two is Multiplied two itself three times similarly we have 2 to the^ 4 where 2 is multiplied by itself four times we have 3 to the^ 2 or also 3 squ where 3 is multiplied by itself twice and we have 4 the^ 1 now any number to the power one basically means that the number is only multiplied by itself a single time which is why any number to the power 1 is the number itself now we are going to understand two main Product of Powers Rule exponent rules the first rule is the product of powers rule this basically says that you have to add the power together when you're multiplying like bases now like bases basically means the same number in the bases let's look at an example here we are going to take 2 to 3 multiplied by 2 ra to 2 now if we expand this and write it simply looks something like 2 2 2 which is basically 2 ^ 3 multiplied by 2 2 which is 2^ 2 now if you notice this is nothing but two multiplied by itself five times so this can be simply written as 2 ra to the^ 5 now this is exactly what the rule tells us when you have like bases with two different exponents and you're multiplying these what you can simply do is you can add the exponents so 2 to the^ 3 2^ 2 simply equals to 2 to^ 3 + 2 which is 2 ra to the ^ 5 like we saw before our next rule is very similar to the previous one except this time Quotient of Power Rule instead of multiplication and addition we have Division and subtraction this rule is called quotient of power rules and it says that you need to subtract the powers when dividing like bases let's take the same example which is 2^ 3 / by 2^ 2 now when you expand it you can simply write this as 2 2 2 / 2 2 now we can simply cancel the twos out and we are left with 2 ided by 1 now 2 ided 1 is obviously only 2 and if you recall any number raised to the power 1 is the number itself so we can write this two as 2 raed to 1 which brings us to the second rule which states that if you have to like bases and they are divided you can simply subtract the exponents so 2 the^ 3 / by 2^ 2 will simply be equal to 2 to 3 - 2 which is 2^ 1 like we saw earlier so now that your Basics are clear let's get into the main question X raed to the^ 0 isal to 1 Main Question now we are going to take the number three let's say that 3 to the^ 0 isal Al to 1 so I'm going to start with the right hand side number which is one and I'm going to prove that this one is equal to 3 ^ 0 now the number one can simply be written as 3 / 3 now if you remember any number raed to the^ 1 is the number itself so we can write this equation as 3^ 1 / 3^ 1 now if you remember your second rule of division we can simply write this equation as 3 raed to 1 - 1 and what is 1 - 1 exactly it's 0 so this is how 1 = to 3 raed to 0 and you can try this with any other number you'll always land on the same solution so I hope you learned something new today if there's any proof or any question related to mathematics like this let me know in the comments below so we can do that and keep following the art of science for |
13771 | https://thirdspacelearning.com/us/math-resources/topic-guides/ratio-and-proportion/percent-decrease/ | Percent Decrease - Math Steps, Examples & Questions
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Introduction
What is percent decrease?
Common Core State Standards
How to decrease a value by a percent
Percent decrease examples
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Example 1: percent decreaseExample 2: percent decreaseExample 3: percent decreaseExample 4: a markdown
Percent decrease using a percent as a decimal
↓
Example 5: percent as a decimalExample 6: percent as a decimal
Calculating percent decrease
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Example 7: calculating percent changeExample 8: calculating percent loss
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Math resourcesRatio and proportionPercent
Percent decrease
Percent decrease
Here you will learn about percent decrease, including how to decrease a value by a given percent, use multipliers to calculate percent decrease, and work out percent change.
Students will first learn about percent decrease as part of ratios and proportions in the 7th grade.
What is percent decrease?
Percent decreasemeans subtracting a given percentage of a value from the original value. To do this, you can either calculate the given percent of the value and subtract it from the original or use a percent as a decimal.
For example,
Decrease \$30 by 20\%.
When given two values, you can calculate the percent difference or percentage change. When the value goes down this may also be called percent decrease, percent loss, or a markdown.
You can calculate percent change using the percentage change formula:
\text { Percent change }=\frac{\text { amount of change }}{\text { original }} \times 100
Percent loss
Percent loss is when the percent decreases or gets smaller between a starting percent and the final percent.
For example,
At the beginning of the month, you had \$150 in your savings account, and by the end of the month, you had \$75. Work out the percent decrease.
\text { Percent decrease }=\cfrac{150-75}{150} \times 100=50 \%
Markdowns
Markdowns are the decreased price of a good. Markdowns are more commonly known as discounts, or sales, in a retail store.
For example, at the end of the summer, stores will place their summer clothing and swimsuits on markdown to make room for their fall and winter clothing.
What is percent decrease?
Common Core State Standards
How does this relate to 7th grade math?
Grade 7: Ratios and Proportional Relationships (7.RP.3)Use proportional relationships to solve multistep ratio and percent problems. Examples: simple interest, tax, markups and markdowns, gratuities and commissions, fees, percent increase and decrease, percent error.
[FREE] Percents Worksheet (Grade 6 to 7)
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[FREE] Percents Worksheet (Grade 6 to 7)
Use this quiz to check your grade 6 to 7 students’ understanding of percents. 10+ questions with answers covering a range of 6th and 7th grade percent topics to identify areas of strength and support!
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How to decrease a value by a percent
In order to decrease a value by a percent:
Calculate the given percent of the value.
Subtract the calculated percent from the original number.
Percent decrease examples
Example 1: percent decrease
Decrease \$80 by 30\%.
Calculate the given percent of the value.
You will need to calculate 30\% of \$80.
The easiest way to calculate 30\% is to calculate 10\% (by dividing by 10 ) and then multiply by 3 to get 30\%.
10\% of \$80 = \$8
8 \times 3=24
30\% of \$80 = \$24
2Subtract the calculated percent from the original number.
You will subtract \$24 from the initial value of \$80.
\$80-\$24 = \$56
A 30\% decrease of \$80 is \$56.
Example 2: percent decrease
Decrease 250 \, km by 45\%.
Calculate the given percent of the value.
You will calculate 45\% of 250 by first finding 40\%.
10\% of 250 \, km = 25 \, km
25 \times 4=100
40\% of 250 \, km = 100 \, km
You will then need to calculate 5\%. Find 10\% then divide by 2 to find 5\%.
10\% of 250 \, km = 25 \, km
25 \div 2=12.5
5\% of 250 \, km = 12.5 \, km
100 \mathrm{~km}+12.5 \mathrm{~km}=112.5 \mathrm{~km}
45\% of 250 \, km = 112.5 \, km
Subtract the calculated percent from the original number.
You will subtract 112.5 \, km from the original value.
250-112.5 = 137.5 \, km
A 45\% decrease of 250 \, km is 137.5 \, km.
Example 3: percent decrease
Decrease 760 by 40\%.
Calculate the given percent of the value.
You will calculate 40\% of 760.
10\% of 760 = 76
76 \times 4=304
40\% of 760 = 304
Subtract the calculated percent from the original number.
You will subtract 304 from the original number.
760-304=456
A 40\% decrease of 760 is 456.
Example 4: a markdown
A jacket that costs \$35.00 is reduced by 15\% in a sale. Calculate the new price of the jacket.
Calculate the given percent of the value.
You will calculate 15\% of \$35.
You will calculate 15\% of \$35 by first finding 10\%.
10\% of \$35 = \$3.50
You will then need to calculate 5\%. Use 10\% then divide by 2 to find 5\%.
10\% of \$35 = \$3.50
3.50 \div 2=1.75
5\% of \$35 = \$1.75
\$ 3.50+1.75=\$ 5.25
15\% of \$35 is \$5.25.
Subtract the calculated percent from the original number.
You will subtract \$5.25 from the original price.
\$ 35.00-\$ 5.25=\$ 29.75
The new price for the jacket is \$29.75.
Percent decrease using a percent as a decimal
In order to decrease a value by using a percent as a decimal:
Subtract the percent you are decreasing from \bf{100\%} .
Convert the percent to a decimal.
Multiply the original amount by the decimal.
Example 5: percent as a decimal
Decrease 84 \, m by 46\%.
Subtract the percent you are decreasing from \bf{100\%} .
Here you are calculating a 46 percent decrease so subtract 46\% from 100\%.
100 \%-46 \%=54 \%
Convert the percent to a decimal.
To convert a percent to a decimal, you divide the percentage by 100.
54 \div 100=0.54
54\% = 0.54
Multiply the original amount by the decimal.
You will multiply the starting value by the decimal.
84\times 0.54=45.36 \, m
84 \, m decreased by 54\% is 45.36 \, m.
Example 6: percent as a decimal
Daniel has \$3,600. He spends 25\% of his money. How much does he have left?
Subtract the percent you are decreasing from \bf{100\%} .
This is a 25 percent decrease so subtract 25\% from 100\%.
100\%-25\% = 75\%
Convert the percentage to a decimal.
To convert a percent to a decimal, you divide the percent by 100.
75 \div 100=0.75
75\% = 0.75
Multiply the original amount by the decimal.
You will multiply \$3,600 by 0.75.
3,600 \times 0.75=\$ 2,700
Daniel has \$2,700 left.
Calculating percent decrease
In order to calculate the percent decrease after a percent change:
Find the amount of change by subtracting the original number from the new number.
Plug numbers into the percent change formula.
\text { Percentage change }=\cfrac{\text { Change }}{\text { Original }} \times 100
Simplify the fraction, if necessary, using equivalent fractions.
Convert the fraction to a decimal.
Calculate percent change.
Example 7: calculating percent change
Ricky weighed 70 \, kg in March. By June his weight had decreased to 62 \, kg. Calculate the percent decrease in his weight.
Find the amount of change by subtracting the original number from the new number.
Ricky’s weight has changed from 70 \, kg to 62 \, kg. You will subtract his original weight from his new weight.
70 \, kg-62 \, kg = 8 \, kg
Plug numbers into the percent change formula.
The change is 8 \, kg and the original amount is 70 \, kg.
\begin{aligned} &\text { Percent change }=\cfrac{\text { Change }}{\text { Original }} \times 100\\ &\text { Percentage change }=\cfrac{8}{70} \times 100 \end{aligned}
Simplify the fraction, if necessary, using equivalent fractions.
\cfrac{8}{70} \, can be simplified to \, \cfrac{4}{35} \, .
\text { Percent change }=\cfrac{4}{35} \times 100
Convert the fraction to a decimal.
\cfrac{4}{35}=0.114
You can round this decimal to 0.11 to solve.
\text { Percent change }=0.11 \times 100
Calculate percent change.
0.11\times 100=11
\text{Percent change } = 11\%
The percent decrease is about 11\%.
Example 8: calculating percent loss
Louise buys a car for \$7,500 and sells it for \$6,150 two years later. Calculate Louise’s percent loss.
Find the amount of change by subtracting the original number from the new number.
The value has changed from \$7,500 to \$6,150.
\$7,500-\$6,150 = \$1,350
Plug numbers into the percent change formula.
The change is \$1,350 and the original amount is \$7,500.
\text { Percent change }=\cfrac{\text { Change }}{\text { Original }} \times 100
\text { Percent change }=\cfrac{1350}{7500} \times 100
Simplify the fraction, if necessary, using equivalent fractions.
\cfrac{1,350}{7,500} \, can be simplified to \, \cfrac{18}{100} \, .
\text { Percent change }=\cfrac{18}{100} \times 100
Convert the fraction to a decimal.
\cfrac{18}{100}=0.18
\text { Percent change }=0.18 \times 100
Calculate percent change.
0.18 \times 100=18
\text{Percent change }=18 \%
Louise’s percent loss is 18\%.
Teaching tips for percent decrease
There are times when the figures used to calculate percent change can be rather big. To help them with these tasks, students can use calculators or percentage change calculators.
Microsoft Excel may be used to calculate percent change. This could serve as an additional teaching point when talking about percent change with students who are comfortable with Excel.
Easy mistakes to make
Incorrectly converting percents to decimals
The most common mistakes are with single digit percents (6\%), multiples of 10 \, (60\%) and decimal percentages (1.2\%).
Remember to divide the percent by 100 to find the decimal.
For example, 6\%=0.06, \, 50\%=0.5, \, 1.2\%=0.012, \, 206.5\%=2.065
Using an incorrect value for the denominator in the percentage decrease formula
Using the new value instead of the original value for the denominator when calculating percentage change.
Related lessons on percent
Percent
Percent of a number
Simple interest
Percent change
Percent increase
Percent increase and decrease
Percent error
Exponential decay
Compound interest formula
Practice percent decrease questions
Calculate the percent decrease that has occurred in each of these questions.
Decrease 7,500 \, m by 20\%.
7,480 \, m
7,520 \, m
9,000 \, m
6,000 \, m
First, you will need to calculate 20\% of 7,500.
10\% of 7,500 \, m = 750 \, m
750 \times 2=1,500
20\% of 7,500 \, m = 1,500 \, m
Then you will subtract 1,500 \, m from the original value.
7,500-1,500=6,000 \mathrm{~m}
A 20\% decrease of 7,500 \, m is 6,000 \, m.
Decrease 45 \, ml by 80\%.
9 \, ml
80 \, ml
36 \, ml
51 \, ml
First, you will need to calculate 80\% of 45 \, ml.
10\% of 45 \, ml = 4.5 \, ml
4.5 \times 8=36
80\% of 45 \, ml = 36 \, ml
Then you will subtract 36 \, ml from the original value.
45-36=9 \, ml
An 80\% decrease of 45 \, ml is 9 \, ml.
Use a percent as a decimal to decrease \$750 by 25\%.
\$562.50
\$731
\$173
\$927.42
This is a 25 percent decrease so start by subtracting 25\% from 100 \%.
100 \%-25 \%=75 \%
Next, to convert a percent to a decimal, you divide the percent by 100.
75 \div 100=0.75
75\% = 0.75
Then, you will multiply the original number, \$750, by the decimal.
750 \times 0.75=\$ 562.50
\$750 decreased by 25\% is 562.50.
Use a percent as a decimal to decrease 250 \, g by 40\%.
215.8 \, g
6.65 \, g
97.028 \, g
150 \, g
This is a 40 percent decrease so start by subtracting 40\% from 100\%.
100 \%-40 \%=60 \%
Next, to convert a percent to a decimal, you divide the percent by 100.
60 \div 100=0.60
60\% = 0.60
Then, you will multiply the original number, 250 \, g, by the decimal.
250 \times 0.60=150
250 \, g decreased by 40\% is 150 \, g.
Find the percent decrease when 650 \, kg is decreased to 320 \, kg.
330\%
50.77 \%
49.23\%
203.125\%
Start by finding the amount of change by subtracting the original number from the new number.
The value has changed from 650 \, kg to 320 \, kg.
650-320=330 \, kg
Next, plug numbers into the percent change formula.
The change is 330 \, kg and the original amount is 650 \, kg.
\text { Percent change }=\cfrac{\text { Change }}{\text { Original }} \times 100
\text { Percent change }=\cfrac{330}{650} \times 100
If the fraction can be simplified, do that next.
\cfrac{330}{650} \, can be simplified to \, \cfrac{33}{65} \, .
\text {Percent change }=\cfrac{33}{65} \times 100
Then you will convert the fraction to a decimal.
\cfrac{33}{65}=0.51
\text {Percent change }=0.51 \times 100
Then you will calculate the percent change.
0.51 \times 100=51
\text {Percent change }=51 \%
The percent change is a 51\% decrease.
Find the percent loss when Tristan buys a tractor for \$11,500 and sells the tractor for \$9,500 one year later.
18.25\%
85 \%
17.4\%
82.25\%
Start by finding the amount of change by subtracting the original number from the new number.
The value has changed from \$11,500 to \$9,500.
11,500-9,500=\$ 2,000
Next, plug numbers into the percent change formula.
The change is \$2,000 and the original amount is \$11,500.
\text { Percent change }=\cfrac{\text { Change }}{\text { Original }} \times 100
\text { Percent change }=\cfrac{2,000}{11,500} \times 100
If the fraction can be simplified, do that next.
\cfrac{2,000}{11,500} \, can be simplified to \, \cfrac{4}{23} \, .
\text {Percent change }=\cfrac{4}{23} \times 100
Then you will convert the fraction to a decimal.
\cfrac{4}{23}=0.174
\text {Percent change }=0.174 \times 100
Then you will calculate the percent change.
0.174 \times 100=17.4
\text {Percent change }=17.4 \%
The tractor decreased in value by 17\%.
Percent decrease FAQs
Can finding the percent change give you a negative number for a percent?
If you are finding a percent decrease, or any percentage where the value is decreasing, the initial percent change calculation will result in a negative number. However, you will use the absolute value of the difference to calculate the overall percent change, which will result in a positive number.
Is there a different formula to calculate percent increase and percent decrease?
You use the same formula whether you are finding the percent increase or decrease. With both, you are calculating the percent to which the measure either increases or decreases.
The next lessons are
Compound measures
Exponents
Converting fractions, decimals, and percentages
Algebraic expressions
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Introduction
What is percent decrease?
Common Core State Standards
How to decrease a value by a percent
Percent decrease examples
↓
Example 1: percent decreaseExample 2: percent decreaseExample 3: percent decreaseExample 4: a markdown
Percent decrease using a percent as a decimal
↓
Example 5: percent as a decimalExample 6: percent as a decimal
Calculating percent decrease
↓
Example 7: calculating percent changeExample 8: calculating percent loss
Teaching tips for percent decrease
Easy mistakes to make
Related lessons on percent
Practice percent decrease questions
Percent decrease FAQs
Next lessons
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13772 | https://maxplanckneuroscience.org/ribosomes-on-the-road-neurons-assemble-protein-factories-far-from-the-nucleus/ | Ribosomes On the Road: Neurons Assemble Protein Factories Far from the Nucleus - Max Planck Neuroscience
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Neural Excitability, Synapses, and GliaResearch News
Ribosomes On the Road: Neurons Assemble Protein Factories Far from the Nucleus
September 3, 2025Max Planck Institute for Brain Research
Ribosome-building blocks are found throughout neurons, including at synapses — suggesting local ribosome maturation. © Fusco et al. 2025 (PNAS), CC BY 4.0.
New study shows ribosome biogenesis can occur within neuronal branches, challenging traditional textbook biology
Ribosomes, the cell’s protein-making machines, have long been thought to assemble entirely within and near the nucleus before being transported to other parts of the cell — including the long extensions of neurons. New research from the Schuman Lab at the Max Planck Institute for Brain Research, published in PNAS, now overturns this view. The study shows that neurons can complete the final steps of ribosome assembly locally within their distal processes — near synapses and far from the nucleus.
“We’ve known for some time that neurons synthesize proteins locally at synapses,” says Dr. Claudia Fusco, first author of the study. “What’s remarkable is that they may also assemble the very machines that make those proteins — right there, on-site.”
A decentralized strategy for protein production
Neurons face a unique logistical challenge: their cell bodies may be located hundreds of microns away from the sites where proteins are needed most — the synapses, where electrical and chemical signals are exchanged. To meet this demand, neurons have evolved a strategy of local translation: delivering messenger RNAs (mRNAs) into dendrites and axons, where ribosomes can produce proteins precisely where they’re needed. However, ribosome biogenesis — the intricate, multistep process of constructing ribosomes from ribosomal RNA (rRNA) and proteins — was still thought to be restricted to the nucleus and its immediate surroundings. The new study shows that this is not the case.
Detecting the machinery in the periphery
To explore the possibility of local ribosome assembly, Fusco and colleagues used a suite of high-resolution techniques, including mass spectrometry, RNA sequencing, imaging, and biochemical fractionation. By isolating neurite-only fractions from cultured rat neurons, they were able to identify the molecular components present in distal processes. Surprisingly, they found all components involved in the cytoplasmic stages of ribosome maturation. These findings indicate that neuronal processes contain everything needed to complete the final assembly and maturation of ribosomes — potentially enabling the on-demand production of functional ribosomes near synapses.
“Local ribosome maturation provides a flexible and efficient mechanism to dynamically regulate protein synthesis in response to synaptic activity or environmental changes. Rather than depending solely on ribosomes transported from the soma, neurons may maintain a pool of immature ribosomal particles that can be rapidly matured and activated when needed“, explains Prof. Erin Schuman, Director of the Department for Synaptic Plasticity at the Max Planck Institute for Brain Research.
The study opens new avenues for exploring how local ribosome biogenesis might contribute to synaptic plasticity, memory formation, and even neurological disorders linked to dysregulated protein synthesis.
C.M. Fusco, A. Staab, A.M. Bourke, G. Tushev, K. Desch, E.M. Lins, E. Ciirdaeva, S. tom Dieck, N. Kaltenschnee, A. Heckel, J.D. Langer, & E.M. Schuman, Neuronal processes contain the essential components for the late steps of ribosome biogenesis, Proc. Natl. Acad. Sci. U.S.A. 122 (31) e2502424122, (2025).Link
Max Planck Institute for Brain Research
brain researchlocal protein synthesisRibosome Biogenesissynaptic plasticity
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13773 | https://pubmed.ncbi.nlm.nih.gov/33073648/ | The impact of cryotherapy for symptomatic cervical ectropion on female sexual function and quality of life - PubMed
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J Obstet Gynaecol
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. 2021 Jul;41(5):815-820.
doi: 10.1080/01443615.2020.1803243. Epub 2020 Oct 19.
The impact of cryotherapy for symptomatic cervical ectropion on female sexual function and quality of life
Sukru Yildiz1,Ismail Alay1,Ecem Eren1,Ibrahim Karaca2,Gul Gultekin3,Cihan Kaya1,Huseyin Cengiz4
Affiliations Expand
Affiliations
1 Department of Obstetrics and Gynecology, Bakirkoy Dr. Sadi Konuk Training and Research Hospital, University of Health Sciences, Istanbul, Turkey.
2 Department of Obstetrics and Gynecology, University of Health Sciences Tepecik Training and Research Hospital, Izmir, Turkey.
3 Department of Family Medicine, Bakirkoy Dr. Sadi Konuk Training and Research Hospital, University of Health Sciences, Istanbul, Turkey.
4 Department of Obstetrics and Gynecology, Faculty of Medicine, Istanbul Aydın University, Istanbul, Turkey.
PMID: 33073648
DOI: 10.1080/01443615.2020.1803243
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Observational Study
The impact of cryotherapy for symptomatic cervical ectropion on female sexual function and quality of life
Sukru Yildiz et al. J Obstet Gynaecol.2021 Jul.
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. 2021 Jul;41(5):815-820.
doi: 10.1080/01443615.2020.1803243. Epub 2020 Oct 19.
Authors
Sukru Yildiz1,Ismail Alay1,Ecem Eren1,Ibrahim Karaca2,Gul Gultekin3,Cihan Kaya1,Huseyin Cengiz4
Affiliations
1 Department of Obstetrics and Gynecology, Bakirkoy Dr. Sadi Konuk Training and Research Hospital, University of Health Sciences, Istanbul, Turkey.
2 Department of Obstetrics and Gynecology, University of Health Sciences Tepecik Training and Research Hospital, Izmir, Turkey.
3 Department of Family Medicine, Bakirkoy Dr. Sadi Konuk Training and Research Hospital, University of Health Sciences, Istanbul, Turkey.
4 Department of Obstetrics and Gynecology, Faculty of Medicine, Istanbul Aydın University, Istanbul, Turkey.
PMID: 33073648
DOI: 10.1080/01443615.2020.1803243
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In rare cases, cervical ectropion causes symptoms such as abundant leucorrhoea, postcoital bleeding, recurrent cervicitis, pelvic pain, and dyspareunia. Cryotherapy is a commonly used treatment for symptomatic cervical ectropion. We assessed the impact of cryotherapy on sexual function and quality of life among patients with symptomatic cervical ectropion. In this prospective observational study, 73 patients were assessed before and six months after cryotherapy treatment using the Female Sexual Function Index (FSFI) and Short Form-12 Health Survey questionnaires. The double-freeze cryotherapy procedure was performed using a cryotherapy unit, and liquid nitrogen was used as a refrigerant. The mean physical and mental quality of life scores were significantly improved after treatment. With the exception of the pain domain, the overall and domain FSFI scores exhibited no significant differences before and after cryotherapy. The sexual pain domain scores were significantly increased after treatment. There was a statistically significant improvement in vaginal discharge, pelvic pain, and postcoital bleeding symptoms after the cryotherapy. We concluded that cryotherapy is an effective and feasible treatment for symptomatic cervical ectropion. Although cryotherapy results in improved quality of life scores, it has no significant impact on female sexual function.Impact statementWhat is already known on this subject? Cryotherapy is the most preferred treatment option for symptomatic cervical ectropion. Its feasibility and effectiveness with respect to symptom relief have been observed in previous studies. No study has evaluated quality of life and sexual function after cryotherapy among patients with symptomatic cervical ectropion.What do the results of this study add? Although the patients' quality of life scores were significantly improved after treatment, no significant improvement was observed in overall and domain sexual function scores, with the exception of the pain domain. The sexual pain domain scores were significantly improved after cryotherapy.What are the implications of these findings for clinical practice and/or further research? Patients should not expect better sexual function after cryotherapy. Comparative studies should seek to identify the ideal treatment option, which would result in both symptom relief and better sexual function.
Keywords: Cervical ectropion; cryotherapy; female sexual function; quality of life.
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Female sexual function and activity following cystectomy and continent urinary tract diversion for benign indications: a clinical pilot study and review of literature.Elzevier HW, Nieuwkamer BB, Pelger RC, Lycklama à Nijeholt AA.Elzevier HW, et al.J Sex Med. 2007 Mar;4(2):406-16. doi: 10.1111/j.1743-6109.2006.00257.x.J Sex Med. 2007.PMID: 17367436 Review.
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Description Formaldehyde-2D.svg Deutsch: Struktur von Formaldehyd (Methanal), CH 2 O
English: Structure of formaldehyde (methanal), CH 2 O
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13775 | https://learning.closer.ac.uk/learning-modules/study-design/methods-of-data-collection/data-collection-instruments/ | Glossary
Administrative data
Administrative data is the term used to describe everyday data about individuals collected by government departments and agencies. Examples include exam results, benefit receipt and National Insurance payments.
Age effects
Age effects relates to changes in an outcome as a result of getting older.
Anonymisation
Anonymisation refers to the removal of study participants' identifying information (e.g., name, address) in order to preserve their privacy.
Attrition
Attrition is the discontinued participation of study participants in a longitudinal study. Attrition can reflect a range of factors, from the study participant not being traceable to them choosing not to take part when contacted. Attrition is problematic both because it can lead to bias in the study findings (if the attrition is higher among some groups than others) and because it reduces the size of the sample.
Baseline
Baseline refers to the start of a study when initial information is collected on participation (however, in longitudinal studies, researchers may adopt an alternative ‘baseline’ for the purposes of analysis).
Biological samples
Biological samples is the term used for specimens collected from human subjects from which biological information, such as genetic markers, can be extracted for analysis. Common examples include blood, saliva or hair.
Body mass index
Body mass index is a measure used to assess if an individual is a healthy weight for their height. It is calculated by dividing the individual’s weight by the square of their height, and it is typically represented in units of kg/m2.
Boosted samples
Boosted samples are used to overcome sample bias due to attrition or to supplement the representation of smaller sub-groups within the sample. Inclusion of boosted samples must be accompanied by appropriate survey weights.
CAPI
Computer-assisted personal interviewing (CAPI) is a technique for collecting data from participants using computers to eliminate common errors such as questionnaire routing and data entry mistakes. The use of computers take place within the context of a face-to-face interview.
CASI
Computer-assisted self-interviewing (CASI) is a technique for collecting data from participants using computers to eliminate common errors such as questionnaire routing and data entry mistakes. The use of computers take place within the context of a self-completion questionnaire.
Categorical variable
A categorical variable is a variable that can take one of a limited number of discrete values. They can be either nominal - they contain no inherent order of categories (e.g. sex; marital status) - or ordinal - they can be ranked in some meaningful order (e.g. level of satisfaction with a service).
CATI
Computer-assisted telephone interviewing (CATI) is a technique for collecting data from participants using computers to eliminate common errors such as questionnaire routing and data entry mistakes. The use of computers take place within the context of a telephone interview.
Censoring
For some study participants the exact time of an event will not be known because either: the study ends (or the analysis is carried out) before they have had the event, or the participant drops out of the study before experiencing the event. It is therefore, only known that the event has not occurred up to the time that they were last observed in the study.
Census
Census refers to a universal and systematic collection of data from all individuals within a population. In the UK, the government conducts a census every ten years with the next one due in 2021.
Codebook
A codebook is a document (online or hard-copy) that contains all the information about how a dataset has been coded, such that it can be deciphered by a researcher not familiar with the original coding frame.
Coding
Coding is the process of converting survey responses into numerical codes to facilitate data analysis. All potential responses (as well as possible reasons for non-response) for each variable are assigned numerical values according to a coding frame.
Cognitive assessments
Cognitive assessments are exercises used to measure thinking abilities, such as memory, reasoning and language. Longitudinal studies collecting data in this way can track the extent to which someone's cognitive abilities change (develop or decline) over time.
Cohort studies
Cohort studies are concerned with charting the lives of groups of individuals who experience the same life events within a given time period. The best known examples are birth cohort studies, which follow a group of people born in a particular period.
Complete case analysis
Complete case analysis is the term used to describe a statistical analysis that only includes participants for which we have no missing data on the variables of interest. Participants with any missing data are excluded.
Conditioning
Conditioning refers to the process whereby participants' answers to some questions may be influenced by their participation in the study – in other words, their responses are 'conditioned' by their being members of a longitudinal study. Examples would include study respondents answering questions differently or even behaving differently as a result of their participation in the study.
Confounding
Confounding occurs where the relationship between independent and dependent variables is distorted by one or more additional, and sometimes unmeasured, variables. A confounding variable must be associated with both the independent and dependent variables but must not be an intermediate step in the relationship between the two (i.e. not on the causal pathway). For example, we know that physical exercise (an independent variable) can reduce a person’s risk of cardiovascular disease (a dependent variable). We can say that age is a confounder of that relationship as it is associated with, but not caused by, physical activity and is also associated with coronary health. See also ‘unobserved heterogeneity’, below.
Continuous variable
A continuous variable is a variable that has an infinite number of uncountable values e.g. time, temperature. They are also known as quantitative variables or scale variables.
Cohort effects
Cohort effects relates to changes in an outcome associated with being a member of a specific cohort of people (e.g. born in the same year; or starting school at the same time).
Coverage
In metadata management, coverage refers to the temporal, spatial and topical aspects of the data collection to describe the comprehensiveness of a dataset. For longitudinal studies, this can relate to the topics that are covered across waves, the population to which one can generalise or the geographic extent of the dataset.
Cross-sectional
Cross-sectional surveys involve interviewing a fresh sample of people each time they are carried out. Some cross-sectional studies are repeated regularly and can include a large number of repeat questions (questions asked on each survey round).
Data access agreement
Within the context of data protection, a data access agreement specifies the terms under which users are provided access to specified datasets. This usually forms part of the application process to the data controller to ensure that researchers adhere to a set of terms regarding data confidentiality, sensitivity and dissemination before accessing the data. See also: research ethics
Data cleaning
Data cleaning is an important preliminary step in the data analysis process and involves preparing a dataset so that it can be correctly analysed. 'Cleaning' the data usually involves identifying data input errors, assessing the completeness of the dataset and verifying any anomalies (e.g. outliers).
Data confidentiality
Within the context of data protection, data confidentiality is the process of protecting participants' data from being accessed or disclosed by those unauthorised to do so. Key methods employed in data confidentiality include anonymisation of responses (removal of personal identifying information) and data encryption (protecting the data using codes and/or passwords).
Data harmonisation
Data harmonisation involves retrospectively adjusting data collected by different surveys to make it possible to compare the data that was collected. This enables researchers to make comparisons both within and across studies. Repeating the same longitudinal analysis across a number of studies allows researchers to test whether results are consistent across studies, or differ in response to changing social conditions.
Data imputation
Data imputation is a technique for replacing missing data with an alternative estimate. There are a number of different approaches, including mean substitution and model-based multivariate approaches.
Data linkage
Data linkage simply means connecting two or more sources of administrative, educational, geographic, health or survey data relating to the same individual for research and statistical purposes. For example, linking housing or income data to exam results data could be used to investigate the impact of socioeconomic factors on educational outcomes.
Data protection
Data protection refers to the broad suite of rules governing the handling and access of information about people. Data protection principles include confidentiality of responses, informed consent of participants and security of data access. These principles are legally protected by the Data Protection Act (DPA) and the General Data Protection Regulation (GDPR).
Data structure
Data structure refers to the way in which data are organised and formatting in advance of data analysis.
Dependent variable
In analysis, the dependent variable is the variable you expect to change in response to different values of your independent (or predictor) variables. For example, a students’ test results may be (partially) explained by the number of hours spent on revision. In this case, the dependent variable is students’ test score, which you expect to be different according to the amount of time spent revising.
Derived variable
A derived variable is a variable that is calculated from the values of other variables and not asked directly of the participants. It can involve a mathematical calculation (e.g. deriving monthly income from annual income by dividing by 12) or a recategorisation of one or more existing variables (e.g. categorising monthly income into £500 bands - £0 to £500, £501 to £1,000, etc.)
Diaries
Diaries are a data collection instrument that is particularly useful in recording information about time use or other regular activity, such as food intake. They have the benefit of collecting data from participants as and when an activity occurs. As such, they can minimise recall bias and provide a more accurate record of activities than a retrospective interview.
Dissemination
Dissemination is the process of sharing information – particularly research findings – to other researchers, stakeholders, policy makers, and practitioners through various avenues and channels, including online, written publications and events. Dissemination is a planned process that involves consideration of target audiences in ways that will facilitate research uptake in decision-making processes and practice.
Dummy variables
Dummy variables, also called indicator variables, are sets of dichotomous (two-category) variables we create to enable subgroup comparisons when we are analysing a categorical variable with three or more categories.
Empirical data
Empirical data refers to data collected through observation or experimentation. Analysis of empirical data can provide evidence for how a theory or assumption works in practice.
Fields
In metadata management, fields are the elements of a database which describes the attributes of items of data.
General ability
General ability is a term used to describe cognitive ability, and is sometimes used as a proxy for intelligent quotient (IQ) scores.
Growth curve modelling
Growth curve modelling is used to analyse trajectories of longitudinal change over time allowing us to model the way participants change over time, and then to explore what characteristics or circumstances influence these patterns of longitudinal change.
Hazard rate
Hazard rate refers to the probability that an event of interest occurs at a given time point, given that it has not occurred before.
Health assessments
Health assessments refers to the assessments carried out on research participants in relation to their physical characteristics or function. These can include measurements of height and weight, blood pressure or lung function.
Heterogeneity
Heterogeneity is a term that refers to differences, most commonly differences in characteristics between study participants or samples. It is the opposite of homogeneity, which is the term used when participants share the same characteristics. Where there are differences between study designs, this is sometimes referred to as methodological heterogeneity. Both participant or methodological differences can cause divergences between the findings of individual studies and if these are greater than chance alone, we call this statistical heterogeneity. See also: unobserved heterogeneity.
Household panel surveys
Household panel surveys collect information about the whole household at each wave of data collection, to allow individuals to be viewed in the context of their overall household. To remain representative of the population of households as a whole, studies will typically have rules governing how new entrants to the household are added to the study.
Incentives and rewards
As a way of encouraging participants to take part in research, they may be offered an incentive or a reward. These may be monetary or, more commonly, non-monetary vouchers or tokens. Incentives are advertised beforehand and can act as an aid to recruitment; rewards are a token of gratitude to the participants for giving their time.
Independent variable
In analysis, an independent variable is any factor that may be associated with an outcome or dependent variable. For example, the number of hours a student spends on revision may influence their test result. In this case, the independent variable, revision time (at least partially) 'explains' the outcome of the test.
Informed consent
A key principle of research ethics, informed consent refers to the process of providing full details of the research to participants so that they are sufficiently able to choose whether or not to consent to taking part.
Kurtosis
Kurtosis is sometimes described as a measure of ‘tailedness’. It is a characteristic of the distribution of observations on a variable and denotes the heaviness of the distribution’s tails. To put it another way, it is a measure of how thin or fat the lower and upper ends of a distribution are.
Life course
A person's life course refers to the experiences and stages an individual passes through during their life. It centres on the individual and emphasises the changing social and contextual processes that influence their life over time.
Longitudinal studies
Longitudinal studies gather data about the same individuals ('study participants') repeatedly over a period of time, in some cases from birth until old age. Many longitudinal studies focus upon individuals, but some look at whole households or organisations.
Metadata
Metadata refers to data about data, which provides the contextual information that allows you to interpret what data mean.
Missing data
Missing data refers to values that are missing and do not appear in a dataset. This may be due to item non-response, participant drop-out (or attrition) or, in longitudinal studies, some data (e.g. date of birth) may be collected only in some waves. Large amounts of missing data can be a problem and lead researchers to make erroneous inferences from their analysis. There are several ways to deal with the issue of missing data, from casewise deletion to complex multiple imputation models.
Multi-level modelling
Multi-level modelling refers to statistical techniques used to analyse data that is structured in a hierarchical or nested way. For example. study participants who belong to the same household, or students who attend the same school may be expected to be more similar to each other than to participants in other households or schools (such as sharing similar contextual influences). This similarity means that the data from participants within these households/schools are not independent. Multi-level models can account for variability at both the individual level and the group (e.g. household or school) level.
Non-response bias
Non-response bias is a type of bias introduced when those who participate in a study differ to those who do not in a way that is not random (for example, if attrition rates are particularly high among certain sub-groups). Non-random attrition over time can mean that the sample no longer remains representative of the original population being studied. Two approaches are typically adopted to deal with this type of missing data: weighting survey responses to re-balance the sample, and imputing values for the missing information.
Observational studies
Observational studies focus on observing the characteristics of a particular sample without attempting to influence any aspects of the participants’ lives. They can be contrasted with experimental studies, which apply a specific ‘treatment’ to some participants in order to understand its effect.
Panel studies
Panel studies follow the same individuals over time. They vary considerably in scope and scale. Examples include online opinion panels and short-term studies whereby people are followed up once or twice after an initial interview.
Peer review
Peer review is a method of quality control in the process of academic publishing, whereby research is appraised (usually anonymously) by one or more independent academic with expertise in the subject.
Period effects
Period effects relate to changes in an outcome associated with living during a particular time, regardless of age or cohort membership (e.g. living through times of war, economic recession or global pandemic).
Piloting
Piloting is the process of testing your research instruments and procedures to identify potential problems or issues before implementing them in the full study. A pilot study is usually conducted on a small subset of eligible participants who are encouraged to provide feedback on the length, comprehensibility and format of the process and to highlight any other potential issues.
Population
Population refers to all the people of interest to the study and to whom the findings will be able to be generalized (e.g. a study looking into rates of recidivism may have a [target] population encompassing everyone with a criminal conviction). Owing to the size of the population, a study will usually select a sample from which to make inferences. See also: sample, representiveness.
Percentiles
A percentile is a measure that allows us to explore the distribution of data on a variable. It denotes the percentage of individuals or observations that fall below a specified value on a variable. The value that splits the number of observations evenly, i.e. 50% of the observations on a variable fall below this value and 50% above, is called the 50th percentile or more commonly, the median.
Primary research
Primary research refers to original research undertaken by researchers collecting new data. It has the benefit that researchers can design the study to answer specific questions and hypotheses rather than relying on data collected for similar but not necessarily identical purposes. See also: secondary research
Prospective study
In prospective studies, individuals are followed over time and data about them is collected as their characteristics or circumstances change.
Qualitative data
Qualitative data are non-numeric - typically textual, audio or visual. Qualitative data are collected through interviews, focus groups or participant observation. Qualitative data are often analysed thematically to identify patterns of behaviour and attitudes that may be highly context-specific.
Quantitative data
Quantitative data can be counted, measured and expressed numerically. They are collected through measurement or by administering structured questionnaires. Quantitative data can be analysed using statistical techniques to test hypotheses and make inferences to a population.
Questionnaires
Questionnaires are research instruments used to elicit information from participants in a structured way. They might be administered by an interviewer (either face-to-face or over the phone), or completed by the participants on their own (either online or using a paper questionnaire. Questions can cover a wide range of topics and often include previously-validated instruments and scales (e.g. the Rosenberg Self-Esteem Scale).
Recall error or bias
Recall error or bias describes the errors that can occur when study participants are asked to recall events or experiences from the past. It can take a number of forms – participants might completely forget something happened, or misremember aspects of it, such as when it happened, how long it lasted, or other details. Certain questions are more susceptible to recall bias than others. For example, it is usually easy for a person to accurately recall the date they got married, but it is much harder to accurately recall how much they earned in a particular job, or how their mood at a particular time.
Record linkage
Record linkage studies involve linking together administrative records (for example, benefit receipts or census records) for the same individuals over time.
Reference group
A reference group is a category on a categorical variable to which we compare other values. It is a term that is commonly used in the context of regression analyses in which categorical variables are being modelled.
Regression analysis
Regression analysis refers to analytical techniques that use a mathematical 'model' to predict values of a dependent variable from values of one or many independent variable(s).
Repeated measures
Repeated measures are measurements of the same variable at multiple time points on the same participants, allowing researchers to study change over time.
Representativeness
Representativeness is the extent to which a sample is representative of the population from which it is selected. Representative samples can be achieved through, for example, random sampling, systematic sampling, stratified sampling or cluster sampling.
Research ethics
Research ethics relates to the fundamental codes of practice associated with conducting research. Ethical issues that need to be considered include providing informed consent to participants, non-disclosure of sensitive information, confidentiality and anonymity safeguarding of vulnerable groups, and respect for participants' well-being. Academic research proposals need be approved by an ethics committee before any actual research (either primary or secondary) can begin.
Research impact
Research impact is the demonstrable contribution that research makes to society and the economy that can be realised through engagement with other researchers and academics, policy makers, stakeholders and members of the general public. It includes influencing policy development, improving practice or service provision, or advancing skills and techniques.
Residuals
Residuals are the difference between your observed values (the constant and predictors in the model) and expected values (the error), i.e. the distance of the actual value from the estimated value on the regression line.
Respondent burden
Respondent burden is a catch all phrase that describes the perceived burden faced by participants as a result of their being involved in a study. It could include time spent taking part in the interview and inconvenience this may cause, as well as any difficulties faced as a result of the content of the interview.
Response rate
Response rate refers to the proportion of participants in the target sample who completed the survey. Longitudinal surveys are designed with the expectation that response rates will decline over time so will typically seek to recruit a large initial sample in order to compensate for likely attrition of participants.
Retrospective study
In retrospective studies, individuals are sampled and information is collected about their past. This might be through interviews in which participants are asked to recall important events, or by identifying relevant administrative data to fill in information on past events and circumstances.
Sample
Sample is a subset of a population that is used to represent the population as a whole. This reflects the fact that it is often not practical or necessary to survey every member of a particular population. In the case of birth cohort studies, the larger ‘population’ from which the sample is drawn comprises those born in a particular period. In the case of a household panel study like Understanding Society, the larger population from which the sample was drawn comprised all residential addresses in the UK.
Sample size
Sample size refers to the number of data units contained within a dataset. It most frequently refers to the number of respondents who took part in your study and for whom there is usable data. However, it could also relate to households, countries or other institutions. The size of a sample, relative to the size of the population, will have consequences for analysis: the larger a sample is, the smaller the margin of error of its estimates, the more reliable the results of the analysis and the greater statistical power of the study.
Sampling frame
A sampling frame is a list of the target population from which potential study participants can be selected.
Scales
Scales are frequently used as part of a research instrument seeking to measure specific concepts in a uniform and replicable way. Typically, they are composed of multiple items that are aggregated into one or more composite scores. Examples of standardised scales include the British Ability Scale (BAS); the Malaise Inventory; and the Rosenberg Self-Esteem Scale.
Scatterplot
A scatterplot is a way of visualising the relationship between two continuous variables by plotting the value of each associated with a single case on a set of X-Y coordinates. For example, students' test scores in English and maths can be represented as point on a graph, with each point representing a single student's English (x-axis) and maths (y-axis) score. Looking at data for many students allows us to build up a visualisation of the relationship between students' scores in maths and English.
Secondary research
Secondary research refers to new research undertaken using data previously collected by others. It has the benefit of being more cost-effective than primary research whilst still providing important insights into research questions under investigation.
Skewness
Skewness is the measure of how assymetrical the distribution of observations are on a variable. If the distribution has a more pronounced/longer tail at the upper end of the distribution (right-hand side), we say that the distribution is negatively skewed. If it is more pronounced/longer at the lower end (left-hand side), we say that it is positively skewed.
Statistical model
A statistical model is a mathematical representation of the relationship between variables.
Statistical software
Statistical software packages are specifically designed to carry out statistical analysis; these can either be open-source (e.g. R) or available through institutional or individual subscription (e.g. SPSS; Stata).
Structured metadata
Structured metadata define the relationship between data items/objects to enable computer systems to understand the contextual meaning of the data. It uses standardised content to facilitate the use of metadata for data discovery and sharing, and the relationship between metadata elements.
Study participants
Study participants are the individuals who are interviewed as part of a longitudinal study.
Survey logic
Also called conditional routing (sometimes called ‘filters’), survey logic refers to the flow that takes respondents through a survey. Respondents may be required to answer some questions only if they had provided a relevant response to a previous question. E.g. Only respondents who are currently at university may be asked to answer a question relating to their degree subject. This is important when considering missing data.
Survey weights
Survey weights can be used to adjust a survey sample so it is representative of the survey population as a whole. They may be used to reduce the impact of attrition on the sample, or to correct for certain groups being over-sampled.
Survival analysis
Survival analysis is an analytical technique that uses time-to-event data to statistically model the probability of experiencing an event by a given time point. For example, time to retirement, disease onset or length of periods of unemployment.
Sweep
The term used to refer to a round of data collection in a particular longitudinal study (for example, the age 7 sweep of the National Child Development Study refers to the data collection that took place in 1965 when the participants were aged 7). Note that the term wave often has the same meaning.
Target population
The population of people that the study team wants to research, and from which a sample will be drawn.
Time to event
Time to event refers to the duration of time (e.g. in hours, days, months, etc.) from a defined baseline to the time of occurrence of an event of interest (e.g. diagnosis of an illness, first re-offence following release from prison). Survival analysis can be used to analyse such data.
Tracing (or tracking)
Tracing (or tracking) describes the process by which study teams attempt to locate participants who have moved from the address at which they were last interviewed.
Unobserved heterogeneity
Unobserved heterogeneity is a term that describes the existence of unmeasured (unobserved) differences between study participants or samples that are associated with the (observed) variables of interest. The existence of unobserved variables means that statistical findings based on the observed data may be incorrect.
User guide
Part of the documentation that is usually provided with statistical datasets, user guides are an invaluable resource for researchers. The guides contain information about the study, including the sample, data collection procedures, and data processing. Use guides may also provide information about how to analyse the data, whether there are missing data due to survey logic, and advice on how to analyse the data such the application of survey weights.
Variables
Variables is the term that tends to be used to describe data items within a dataset. So, for example, a questionnaire might collect information about a participant’s job (its title, whether it involves any supervision, the type of organisation they work for and so on). This information would then be coded using a code-frame and the results made available in the dataset in the form of a variable about occupation. In data analysis variables can be described as ‘dependent’ and ‘independent’, with the dependent variable being a particular outcome of interest (for example, high attainment at school) and the independent variables being the variables that might have a bearing on this outcome (for example, parental education, gender and so on).
Vulnerable groups
Vulnerable groups refers to research participants who may be particularly susceptible to risk or harm as a result of the research process. Different groups might be considered vulnerable in different settings. The term can encompass children and minors, adults with learning difficulties, refugees, the elderly and infirm, economically disadvantaged people, or those in institutional care. Additional consideration and mitigation of potential risk is usually required before research is carried out with vulnerable groups.
Wave
The term used to refer to a round of data collection in a particular longitudinal study (for example, the age 7 wave of the National Child Development Study refers to the data collection that took place in 1965 when the participants were aged 7). Note that the term sweep often has the same meaning.
Learning Hub
Study design
3.1
Data collection instruments
Data collection instruments
Questionnaires
Cognitive assessments
Diaries
Health assessments
Assessments of mental health and wellbeing
Biological samples
Qualitative information
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Data collection instruments
This section looks at the instruments and tools used to capture information about different aspects of participants’ lives.
Study teams choose the best possible survey instrument to collect the information they are looking for. Because longitudinal studies collect a broad range of information at each sweep, they require more than one type of survey instrument each time.
About the respondents
Most data collection instruments are completed by the study participant, but as we learned in the introduction to longitudinal studies module, information is sometimes collected from important people in participants’ lives as well.
For instance, if the study participants are too young to answer questions themselves, the information will be collected from their parents. As the participants get older, the information they provide may be supplemented by information from their parents, siblings, teachers, school nurses, health visitors, partners or children.
Challenges for longitudinal studies
Otis Dudley Duncan is credited with saying “if you want to measure change, don’t change the measure”. This is a particular challenge for longitudinal study teams, which often need to balance the appeal of maintaining consistent measures over time with the need to ensure they are collecting the highest quality and most relevant information possible.
For example, imagine a study team is particularly interested in looking at how blood pressure develops over someone’s lifetime and how it relates to different circumstances. In the early sweeps of data collection, they use the best technology available at the time, but at later sweeps they have the choice to move to a more advanced and accurate blood pressure reader that is now typically used on other studies.
Changing to the latter would improve their data quality and help make it comparable with other studies, but could involve new costs (if they have to purchase the new devices) and will make it harder to compare their new data with those collected earlier using the older devices.
A key way of minimising the impact of any change will be to identify or carry out calibration studies that compare measurements using the different machines.
Back [glossary_exclude]Section 2: Sampling: Boosted samples[/glossary_exclude] Next [glossary_exclude]Questionnaires[/glossary_exclude]
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13776 | https://www.ck12.org/book/cbse_physics_book_class_xii/section/9.0/ | Ray Optics | CK-12 Foundation
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HomePhysicsFlexBooksCK-12 CBSE Physics Class 12Ch9
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9.0 Ray Optics
Difficulty Level: Basic | Created by: CK-12
Last Modified: Feb 18, 2016
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9.1 Laws of Reflection
9.2 Reflection of Light by Spherical Mirrors
9.3 Refraction of Light
9.4 Total Internal Reflection
9.5 Refraction at Spherical Surfaces and by Lens
9.6 Refraction Through a Prism
9.7 Dispersion Through a Prism
9.8 Optical Instruments
9.9 Self Assessments
9.10 Home Assignments
9.11 Class Activities
9.12 Reference for Basic Understanding
9.13 Reference for Advanced Understanding
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This physics text was created using CK-12 resources to be seed content for a complete Physics Class 12 course for CBSE students.
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Yashu Chhabra, PGT Physics,Kulachi Hansraj Model School,Delhi
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13777 | https://qknowbooks.gitbooks.io/jhs_2_science-chemical-compounds/naming_of_oxo-anions.html | Naming of Oxo-anions · Chemical Compounds
Chemical Compounds
Introduction
Chemical Compounds
Chemical Formulae
Systematic Naming of Compounds
Naming of Oxo-anions
Naming of Oxo-acids
Chemical Equations
How to Balance Chemical Equations
Summary
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Naming of Oxo-anions
Naming of Oxo-anions
Oxo-anions are two components coming together. They are Oxo from the word oxygen and anions which is a negatively charged atom. Examples are PO 4 3- , MnO 4–.
Rules for naming Oxo-anions.
Find the number of Oxygen atoms in the oxo-anions using the following prefixes: Mono (1 oxygen), di (2 oxygen atoms), tri (3 oxygen atoms), tetra (4 oxygen atoms), pent (5 oxygen atoms), heria (6 oxygen atoms), hepta (7 oxygen atoms), etc.
The name of the oxo-anion which usually ends in ‘ate’ is written. Example, the oxo-anion MnO 4– is known as Manganate.
The oxidation number of the first element of the oxo-anion is calculated and written in a bracket.
Example: MnO 4–
The oxygen atom in the oxo-anion is four (4) and this is due to the 4 in front of the first oxygen (O). The four (4) oxygen atoms in the oxo-anion gives the name tetraoxo (tetra from four (4) and oxo from oxygen). The oxo-anion is manganate.
From the calculation of oxidation number in the previous lessons, the oxidation number of oxygen (O) is -2.
Let’s calculate the number to be written in front of the tetraoxomanganate or the oxo-anion. MnO4 –
The oxidation number of Mn is not known so it will be represented by x. .: x + 4 (-2) = -1
The -1 is from the charge on the oxo-anion which is the (-) on top of the oxygen.
x +4 (-2) = -1
x – 8 = -1
x = -1+8
x = +7
.: the number in front of the tetraoxomanganate will be +7 which will be written in Roman numerals as (VII).
Hence the systematic name of MnO 4– is tetraoxomanganate (VII). The word ion must be written at the end.
.: MnO 4– – tetraoxomanganate ion.
NB: When ion is not written when naming an oxo-anion, the systematic name is very wrong.
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13778 | http://sineof1.github.io/math_notebook/picks_theorem | Pick's Theorem
So, Pick's Theorem involves some interesting terminology, like lattice polygon, that you don't hear in your typical geometry class. And it
can involve a neat methodology, mathematical induction, which isn't talked about enough, if at all, in high school--and which I won't talk about in this installment.
At the moment, you can click and drag on the dot grid above to draw stuff and see the boundary points of your polygons circled, and in the future, I'll hook this up with some more magic to highlight what Pick's theorem talks about. Play with it to see what makes it tick.
So, what does this theorem say? It says that when you draw any lattice
polygon, you can count the points on the interior, add that to half the number of points that lie on the sides (including the vertices), and
subtract 1, and that will give you the area value of the polygon.
As that definition at the link above will tell you, a lattice polygon is a polygon whose vertices are all on the points of a lattice, which is a
"regularly spaced array of points." The figure at the right shows a trapezoid I drew on my homemade lattice above. Pick's Theorem says that the area of
the trapezoid is equal to Ip+12Bp−1. That is, the number of interior points, Ip, plus half
the number of boundary points (including the vertices), Bp, minus 1 is equal to the area.
Compare the result of this formula to what you get when you find the area of the trapezoid using base and height measures: A=12h⋅(b1+b2). And of course, draw your own to see. Think ridiculously simple, like a unit square, for example. Or ridiculously
complex. Try to break the formula.
Always Start Simple
As far as proving stuff, do what mathematicians do--instead of trying to tackle a really abstract and complex problem all at once, start with
simpler instances of the general case and don't tell anyone until you've solved the harder problem. So, instead of "any lattice polygon," start with
"any lattice rectangle."
How many interior points are there in any lattice rectangle? Well, if you draw a rectangle on the lattice above, you will see that the interior
points of a lattice rectangle form a smaller rectangle of their own. If the height of the original rectangle is h units, then the
smaller rectangle will have h−1 points on its height. The same goes for the width, so the number of points on the interior of the
rectangle will be (h−1)(w−1). (You just have to get used to going back and forth between "units" and "points" in this
explanation for it to make the most sense.)
What about boundary points? We can traverse the boundary of a rectangle, again relating the number of points we count to the width and height. So,
going clockwise from the top of the rectangle at the left, first I count w+1 points on the width (its width is 10 units, but there
are 11 points). Then I count h points along the right side (the height of the rectangle is 6 units, and this is the number I count,
because I have already counted the top right vertex point). And so on. Adding up all the points and simplifying, we get 2w+2h for
the number of boundary points of a lattice rectangle.
Finally, Pick's theorem says that
A=2w+2h2+(h−1)(w−1)−1
And this simplifies to hw, the area of the rectangle--and, more generally, any rectangle.
Let's remember to talk more about this whole theorem at some point in the future. In the meantime, here's a great paper on Pick's
Theorem with some nice explanations.
Postscript
To my point above about mathematicians' silence on interim results and process, I came across this wonderful quotation by mathematician Alexander Grothendieck, who passed away in November of
2014. Here he talks about 'speculation' in mathematics (in a 1983 letter):
It is the kind of thing which has traditionally been lacking in mathematics since the very beginnings, I feel, which is one big
drawback in comparison to all other sciences, as far as I know. Of course, no creative mathematician can afford not to "speculate," namely to do more
or
less daring guesswork as an indispensable source of inspiration. The trouble is that, in obedience to a stern tradition, almost nothing of this appears
in writing, and preciously little even in oral communication. |
13779 | https://www.ahschools.us/cms/lib08/MN01909485/Centricity/Domain/772/MCA_ACT_Statistics_Packet.pdf | MCA/ACT Math Additional Practice Name:___ Statistics 1. Cynthia and her father planted a tree in their front yard 8 years ago. The tree was 2 meters in height when it was planted. The scatterplot below shows how the height of the tree increased each year. Which of the following most closely approximates the equation of the line of best fit for the data? A. y = -2x + 2 B. y = 2x + 2 C. y = - ½x + 2 D. y = ½x + 2 2. The equation of the regression line for a student’s grade y given the number of times they are tardy x is y = -5.1x +97.2. What would someone’s grade be who was tardy 7 times? 3. A study of child development measures the age (in months) at which a child begins to talk and also the child’s score on an ability test given several years later. The study asks whether the age at which a child talks helps predict the later test score. The least-squares regression line of test score y on age x is y = 110 – 1.3x. According to this regression line, what happens (on the average) when a child starts talking one month later? (a) The test score goes down 110 points. (b) The test score goes down 1.3 points. (c) The test score goes up 110 points. (d) The test score goes up 1.3 points. (e) The test score is 108.7. 4. Use the histogram to the right to answer the following questions. a. How many employees earned between$4,290 and $13,050? b. In what class does the median most likely fall? 5. The pie graph below shows the expenses of the Taylors family in a particular month. If the Taylors spent $600 on food, how much did they spend that month on clothing? 6. The box and whisker plot shown below represents a class’s scores on a statistics quiz. Which of the following statements is false? a. Half of the students scored between 12 and 21. b. The median score was 15. c. A score of 21 is at the 75th percentile. d. More students scored between 15 and 21 than any other quartile. e. 25% of the students scored below 12. Food 30% Clothing Mortgage 35% 7. What is the interquartile range of the data shown in the boxplot? 8. Which of the following scatterplots shown below would be best represented by a line of best fit with the equation: y = 25 9. Find the 5 number summary (minimum, lower quartile, median, upper quartile, maximum) of the following data: 10. The heights (in inches) of the members of a soccer team are listed below. 66, 61, 71, 62, 64, 70, 64, 63, 72, 68 After a new member joined the team, the median height of all members was 66 inches. Which of the following could be the height of the new member? A. 68 inches B. 65 inches C. 64 inches D. 61 inches 11. Find x and y so that the ordered data set has a mean of 42 and a median of 35. 17, 22, 26, 29, 34, x, 42, 67, 70, y 12. Which of the following statements, related to the mean, mode and median of the data set given below is true? 12, 15, 10, 19, 5, 5 A. The mean is equal to the median B. The mode is larger than the mean C. The median is smaller than the mode D. The above data set has two modes E. The median of the above data is equal to 14.5 13. Tyrone calculated the mean, median, mode and range of the following data set: 1, 2, 5, 5, 5, 8, 11, 13 Then he realized that a data point of “2” was missing from the original data set. The corrected data set is shown below. 1, 2, 2, 5, 5, 5, 8, 11, 13 When Tyrone accurately recalculates the statistics, which measure will change? A. mean B. median C. mode D. range 14. Jason and Eric discovered that the means of their grades for the first marking period in their math class were identical. They also noticed that the standard deviation of Jason’s grades is 20.7 while the standard deviation of Eric’s grades is 2.7. Which statement must be true? A. In general, Eric’s grades were lower than Jason’s grades. B. Eric’s grades are more consistent than Jason’s grades. C. Eric had more failing grades during the marking period than Jason had. D. The median for Eric’s grades is lower than the median for Jason’s grades. 15. Which of the following data sets has the largest standard deviation? A. {1, 2, 3, 4, 5} B. {2,3, 3, 3, 4} C. {2, 2, 2, 4, 5} D. {0, 2, 3, 4, 6} E. {-1, 1, 3, 5, 7} 16. The number of children of each of the first 41 United States presidents is given in the accompanying table. For this population, determine the mean and the standard deviation to the nearest tenth. 17. The following table gives the 2011 regular season win-loss record of all the NFC teams in the National Football League. a. Find the percentile of the number of wins for the Minnesota Vikings. b. What team was at the 81st percentile? 18. A study is conducted to determine whether office workers have high blood pressure. The participants in the study were friends of the researcher who shared the same doctor. Is this study biased? 19. A tally was made of the number of times each color of crayon was used by a kindergarten class. Which measure of central tendency should the teacher use to determine which color is the favorite color of her class? a. mean b. mode c. median d. range 20. Which situation should be analyzed using bivariate data? a. Ms. Saleem keeps a list of the amount of time her daughter spends on her social studies homework. b. Mr. Benjamin tries to see if his students’ shoe sizes are directly related to their heights. c. Mr. DeStefan records his customers’ best video game scores during the summer. d. Mr. Chan keeps track of his daughter’s algebra grades for the quarter. |
13780 | https://danielpollithy.github.io/general/jensens-inequality | Home
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Jensen's inequality
by Daniel Pollithy —
This is an illustrative explanation of Jensen’s inequality applied to probability.
Jensen’s inequality
Remember that a convex function is a function whose area above the graph (called epi graph) is a convex set. Informally speaking, a convex function can be identified by having positive curvature everywhere. A concave function has the opposite attributes.
Jensen’s inequality states for (\theta \in [0,1]) and (f) convex:
[f(\theta x_1 + (1 - \theta) x_2) \le \theta f(x_1) + (1 - \theta) f(x_2)]
The following image illustrates the inequality:
This is usable to define and show convexity of functions analogously to the definition of convex sets.
In probability
If we have a random variable X with probability distribution p(x) over x, then we could set (x_1=x_2=E[X]). For this case the Jensen’s inequality displays itself as:
[f(E[X]) \le E[f(X)]]
For affine functions f the equality holds. Which tells us that they are concave and convex.
Derivation of the ELBO
(log p(X)) is the log evidence, X the observed variables and Z the latent variables. q denotes a variational approximation of p(Z|X):
[log p(X) = log \int_{Z} p(X,Z)] [= log \int_{Z} p(X,Z) \frac{q(Z)}{q(Z)} = log \int_{Z} \frac{p(X,Z)}{q(Z)} q(Z)]
This is the definition of the expected value wrt q(Z):
[= log(E[\frac{p(X,Z)}{q(Z)}])]
Here comes the inverted Jensen’s inequality in to play (it holds for concave functions). We can confirm ourselves easily that the log() is a concave function by noting that the area under its curve is convex.
[log(E[\frac{p(X,Z)}{q(Z)}]) \ge E[log( \frac{p(X,Z)}{q(Z)} )]]
This can be brought into a nicer form by realizing that the definition of entropy is (H(p) = E_{x\sim p(x)}[log(\frac{1}{p(x)})] = -E_{x\sim p(x)}[log ~ p(X)])
[E[log( \frac{p(X,Z)}{q(Z)} )] = E[log(p(X,Z)) - log(q(Z))]] [E[log(p(X,Z)) - log(q(Z))] = E[log(p(X,Z))] - E[log(q(Z))]] [E[log(p(X,Z))] - E[log(q(Z))] = E[log(p(X,Z))] + H(q(Z))]
Proof Kullback-Leibler divergence is always >= 0
The KL between p(x) and q(x) is defined as:
[KL(p(x);q(x)) = -E_{x \sim p(x)}[log(\frac{q(x)}{p(x)})]]
We can apply the concave version of Jensen’s inequality. q/p is the distribution and log the concave function.
[-E_{x \sim p(x)}[log(\frac{q(x)}{p(x)})] \le - log(E_{x \sim p(x)}[\frac{q(x)}{p(x)}])]
Multiply by (-1):
[E_{x \sim p(x)}[log(\frac{q(x)}{p(x)})] \ge log(E_{x \sim p(x)}[\frac{q(x)}{p(x)}])] [log(E_{x \sim p(x)}[\frac{q(x)}{p(x)}]) = log(\int p(x)\frac{q(x)}{p(x)} ~ dx) = log(\int q(x) ~ dx)]
q is a pdf therefore this integral has to be 1. log(1) = 0. As a result:
[E_{x \sim p(x)}[log(\frac{q(x)}{p(x)})] \ge 0]
Author
Daniel Pollithy
daniel@pollithy.com
Developer based in Munich, Germany |
13781 | https://en.wikipedia.org/wiki/Transcription-translation_coupling | Published Time: 2020-09-04T08:05:37Z
Transcription-translation coupling - Wikipedia
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1 Significance
2 MechanismsToggle Mechanisms subsection
2.1 Expressome complex
2.2 Ribosome-mediated attenuation
2.3 Polarity
3 Discovery
4 References
Transcription-translation coupling
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From Wikipedia, the free encyclopedia
Mechanism of gene expression regulation
Transcription-translation coupling is a mechanism of gene expression regulation in which synthesis of an mRNA (transcription) is affected by its concurrent decoding (translation). In prokaryotes, mRNAs are translated while they are transcribed. This allows communication between RNA polymerase, the multisubunit enzyme that catalyzes transcription, and the ribosome, which catalyzes translation. Coupling involves both direct physical interactions between RNA polymerase and the ribosome ("expressome" complexes), as well as ribosome-induced changes to the structure and accessibility of the intervening mRNA that affect transcription ("attenuation" and "polarity").
Significance
[edit]
Bacteria depend on transcription-translation coupling for genome integrity, termination of transcription and control of mRNA stability. Consequently, artificial disruption of transcription-translation coupling impairs the fitness of bacteria. Without coupling, genome integrity is compromised as stalled transcription complexes interfere with DNA replication and induce DNA breaks. Lack of coupling produces premature transcription termination, likely due to increased binding of termination factor Rho. Degradation of prokaryotic mRNAs is accelerated by loss of coupled translation due to increased availability of target sites of RNase E. It has also been suggested that coupling of transcription with translation is an important mechanism of preventing formation of deleterious R-loops. While transcription-translation coupling is likely prevalent across prokaryotic organisms, not all species are dependent on it. Unlike Escherichia coli, in Bacillus subtilis transcription significantly outpaces translation, and coupling consequently does not occur.
Mechanisms
[edit]
Translation promotes transcription elongation and regulates transcription termination. Functional coupling between transcription and translation is caused by direct physical interactions between the ribosome and RNA polymerase ("expressome complex"), ribosome-dependent changes to nascent mRNA secondary structure which affect RNA polymerase activity (e.g. "attenuation"), and ribosome-dependent changes to nascent mRNA availability to transcription termination factor Rho ("polarity").
Expressome complex
[edit]
The expressome is a supramolecular complex consisting of RNA polymerase and a trailing ribosome linked by a shared mRNA transcript. It is supported by the transcription factors NusG and NusA, which interact with both RNA polymerase and the ribosome to couple the complexes together. When coupled by transcription factor NusG, the ribosome binds newly synthesized mRNA and prevents formation of secondary structures that inhibit transcription. Formation of an expressome complex also aids transcription elongation by the trailing ribosome opposing back-tracking of RNA polymerase. Three-dimensional models of ribosome-RNA polymerase expressome complexes have been determined by cryo-electron microscopy.
Ribosome-mediated attenuation
[edit]
Ribosome-mediated attenuation is a gene expression mechanism in which a transcriptional termination signal is regulated by translation. Attenuation occurs at the start of some prokaryotic operons at sequences called "attenuators", which have been identified in operons encoding amino acid biosynthesis enzymes, pyrimidine biosynthesis enzymes and antibiotic resistance factors. The attenuator functions via a set of mRNA sequence elements that coordinate the status of translation to a transcription termination signal:
A short open reading frame encoding a "leader peptide"
A transcription pause sequence
A "control region"
A transcription termination signal
Once the start of the leader open reading frame has been transcribed, RNA polymerase pauses due to folding of the nascent mRNA. This programmed arrest of transcription gives time for translation of the leader peptide to commence, and transcription to resume once coupled to translation. The downstream "control region" then modulates the elongation rate of either the ribosome or RNA polymerase. The factor determining this depends on the function of the downstream genes (e.g. the operon encoding enzymes involved in the synthesis of histidine contains a series of histidine codons is the control region). The role of the control region is to modulate whether transcription remains coupled to translation depending on the cellular state (e.g. a low availability of histidine slows translation leading to uncoupling, while high availability of histidine permits efficient translation and maintains coupling). Finally, the transcription terminator sequence is transcribed. Whether transcription is coupled to translation determines whether this stops transcription. The terminator requires folding of the mRNA, and by unwinding mRNA structures the ribosome elects the formation of either of two alternative structures: the terminator, or a competing fold termed the "antiterminator".
For amino acid biosynthesis operons, these allow the gene expression machinery to sense the abundance of the amino acid produced by the encoded enzymes, and adjust the level of downstream gene expression accordingly: transcription occurring only if the amino acid abundance is low and the demand for the enzymes is therefore high. Examples include the histidine (his) and tryptophan (trp) biosynthetic operons.
The term "attenuation" was introduced to describe the his operon. While it is typically used to describe biosynthesis operons of amino acids and other metabolites, programmed transcription termination that does not occur at the end of a gene was first identified in λ phage. The discovery of attenuation was significant as it represented a regulatory mechanism distinct from repression. The trp operon is regulated by both attenuation and repression, and was the first evidence that gene expression regulation mechanisms can be overlapping or redundant.
Polarity
[edit]
"Polarity" is a gene expression mechanism in which transcription terminates prematurely due to a loss of coupling between transcription and translation. Transcription outpaces translation when the ribosome pauses[citation needed] or encounters a premature stop codon. This allows the transcription termination factor Rho to bind the mRNA and terminate mRNA synthesis. Consequently, genes that are downstream in the operon are not transcribed, and therefore not expressed. Polarity serves as mRNA quality control, allowing unused transcripts to be terminated prematurely, rather than synthesized and degraded.
The term "polarity" was introduced to describe the observation that the order of genes within an operon is important: a nonsense mutation within an upstream gene effects the transcription of downstream genes. Furthermore, the position of the nonsense mutation within the upstream gene modulates the "degree of polarity", with nonsense mutations at the start of the upstream genes exerting stronger polarity (more reduced transcription) on downstream genes.
Unlike the mechanism of attenuation, which involves intrinsic termination of transcription at well-defined programmed sites, polarity is Rho-dependent and termination occurs at variable position.
Discovery
[edit]
The potential for transcription and translation to regulate each other was recognized by the team of Marshall Nirenberg, who discovered that the processes are physically connected through the formation of a DNA-ribosome complex. As part of the efforts of Nirenberg's group to determine the genetic code that underlies protein synthesis, they pioneered the use of cell-free in vitro protein synthesis reactions. Analysis of these reactions revealed that protein synthesis is mRNA-dependent, and that the sequence of the mRNA strictly defines the sequence of the protein product. For this work in breaking in the genetic code, Nirenberg was jointly awarded the Nobel Prize in Physiology or Medicine in 1968. Having established that transcription and translation are linked biochemically (translation depends on the product of transcription), an outstanding question remained whether they were linked physically - whether the newly synthesized mRNA released from the DNA before it is translated, or if can translation occur concurrently with transcription. Electron micrographs of stained cell-free protein synthesis reactions revealed branched assemblies in which strings of ribosomes are linked to a central DNA fibre. DNA isolated from bacterial cells co-sediment with ribosomes, further supporting the conclusion that transcription and translation occur together. Direct contact between ribosomes and RNA polymerase are observable within these early micrographs. The potential for simultaneous regulation of transcription and translation at this junction was noted in Nirenberg's work as early as 1964.
References
[edit]
^ Artsimovitch, I. (2018). "Rebuilding the bridge between transcription and translation". Molecular Microbiology. 108 (5): 467–472. doi:10.1111/mmi.13964. PMC 5980768. PMID 29608805.
^ McGary K. & Nudler, E. (2013). "RNA polymerase and the ribosome: the close relationship". Current Opinion in Microbiology. 16 (2): 112–117. doi:10.1016/j.mib.2013.01.010. PMC 4066815. PMID 23433801.
^ Jump up to: a b Klaholz, B. (2017). "The ribosome holds the RNA polymerase on track in bacteria". Trends in Biochemical Sciences. 42 (9): 686–689. doi:10.1016/j.tibs.2017.07.003. PMID 28801047.
^ Dutta, D.; Shatalin, K.; Epshtein, V.; Gottesman, M. E. & Nudler, E. (2011). "Linking RNA polymerase backtracking to genome instability in E. coli". Cell. 146 (4): 533–543. doi:10.1016/j.cell.2011.07.034. PMC 3160732. PMID 21854980.
^ Zhu, M.; Mori, M.; Hwa, T. & Dai, X. (2019). "Disruption of transcription-translation coordination in Escherichia coli leads to premature transcriptional termination". Nature Microbiology. 4 (12): 2347–2356. doi:10.1038/s41564-019-0543-1. PMC 6903697. PMID 31451774.
^ Iost, I. & Dreyfus, M. (1995). "The stability of Escherichia coli lacZ mRNA depends upon the simultaneity of its synthesis and translation". EMBO Journal. 14 (13): 3252–61. doi:10.1002/j.1460-2075.1995.tb07328.x. PMC 394387. PMID 7542588.
^ Gowrishankar, J. & Harinarayanan, R. (2004). "Why is transcription coupled to translation in bacteria?". Molecular Microbiology. 54 (3): 598–603. doi:10.1111/j.1365-2958.2004.04289.x. PMID 15491353.
^ Johnson, G. E.; Lalanne, J.; Peters, M. L. & Li, G. (2020). "Functionally uncoupled transcription-translation in Bacillus subtilis". Nature. 585 (7823): 124–128. Bibcode:2020Natur.585..124J. doi:10.1038/s41586-020-2638-5. PMC 7483943. PMID 32848247.
^ Jump up to: a b c Webster, M. W.; Takacs, M.; Zhu, C.; Vidmar, V.; Eduljee, A.; Abdelkareem, M. & Weixlbaumer, A. (2020). "Structural basis of transcription-translation coupling and collision in bacteria". Science. 369 (6509): 1355–1359. Bibcode:2020Sci...369.1355W. doi:10.1126/science.abb5036. PMID 32820062. S2CID 221222557.
^ Jump up to: a b O'Reilly, F. J.; Xue, L.; Graziadei, A.; Sinn, L.; Lenz, S.; Tegunov, D.; Blötz, C.; Singh, N.; Hagen, W.; Cramer, P.; Stülke, J.; Mahamid, J. & Rappsilber, J. (2020). "In-cell architecture of an actively transcribing-translating expressome". Science. 369 (6503): 554–557. Bibcode:2020Sci...369..554O. doi:10.1126/science.abb3758. hdl:21.11116/0000-0006-D30E-D. PMC 7115962. PMID 32732422.
^ Jump up to: a b Wang, C.; Molodtsov, V.; Firlar, E.; Kaelber, J.; Blaha, G.; Su, M. & Ebright, R. H. (2020). "Structural basis of transcription-translation coupling". Science. 369 (6509): 1359–1365. Bibcode:2020Sci...369.1359W. doi:10.1126/science.abb5317. PMC 7566311. PMID 32820061. S2CID 221220008.
^ Proshkin, S.; Rahmouni, A. R.; Mironov, A. & Nudler, E. (2010). "Cooperation between translating ribosomes and RNA polymerase in transcription elongation". Science. 328 (5977): 504–508. Bibcode:2010Sci...328..504P. doi:10.1126/science.1184939. PMC 2930199. PMID 20413502.
^ Stevenson-Jones, F.; Woodgate, J.; Castro-Roa, D. & Zenkin, N. (2020). "Ribosome reactivates transcription by physically pushing RNA polymerase out of transcription arrest". Proceedings of the National Academy of Sciences. 117 (15): 8462–8467. Bibcode:2020PNAS..117.8462S. doi:10.1073/pnas.1919985117. PMC 7165469. PMID 32238560.
^ Kohler, R.; Mooney, R. A.; Mills, D. J.; Landick, R. & Cramer, P. (2017). "Architecture of a transcribing-translating expressome". Science. 356 (6334): 194–197. Bibcode:2017Sci...356..194K. doi:10.1126/science.aal3059. PMC 5528865. PMID 28408604.
^ Turnbough, C. L. (2019). "Regulation of Bacterial Gene Expression by Transcription Attenuation". Microbiology and Molecular Biology Reviews. 83 (3). doi:10.1128/MMBR.00019-19. PMC 6710462. PMID 31270135.
^ Yanofsky C (1981). "Attenuation in the control of expression of bacterial operons". Nature. 289 (5800): 751–758. Bibcode:1981Natur.289..751Y. doi:10.1038/289751a0. PMID 7007895. S2CID 4364204.
^ Jump up to: a b Yanofsky C (2000). "Transcription attenuation: once viewed as a novel regulatory strategy". Journal of Bacteriology. 182 (1): 1–8. doi:10.1128/jb.182.1.1-8.2000. PMC 94232. PMID 10613855.
^ Jump up to: a b Kasai, T. (1974). "Regulation of the expression of the histidine operon in Salmonella typhimurium". Nature. 249 (5457): 523–527. Bibcode:1974Natur.249..523K. doi:10.1038/249523a0. PMID 4599761. S2CID 472218.
^ Johnston, H. M.; Barnes, W. M.; Chumley, F. G.; Bossi, L. & Roth, J. R. (1980). "Model for regulation of the histidine operon of Salmonella". Proceedings of the National Academy of Sciences. 77 (1): 508–512. Bibcode:1980PNAS...77..508J. doi:10.1073/pnas.77.1.508. PMC 348301. PMID 6987654.
^ Landick, R.; Carey, J. & Yanofsky, C. (1985). "Translation activates the paused transcription complex and restores transcription of the trp operon leader region". Proceedings of the National Academy of Sciences. 82 (14): 4663–4667. Bibcode:1985PNAS...82.4663L. doi:10.1073/pnas.82.14.4663. PMC 390446. PMID 2991886.
^ Luzzati, D. (1970). "Regulation of lambda exonuclease synthesis: role of the N gene product and lambda repressor". Journal of Molecular Biology. 49 (2): 515–519. doi:10.1016/0022-2836(70)90261-5. PMID 4915096.
^ Singer, C. E.; Smith, G. R.; Cortese, R. & Ames, B. N. (1972). "Mutant tRNA His ineffective in repression and lacking two pseudouridine modifications". Nature New Biology. 238 (81): 72–74. doi:10.1038/newbio238072a0. PMID 4558263.
^ Jackson, E. N. & Yanofsky, C. (1973). "The region between the operator and first structural gene of the tryptophan operon of Escherichia coli may have a regulatory function". Journal of Molecular Biology. 76 (1): 89–101. doi:10.1016/0022-2836(73)90082-x. PMID 4578102.
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13782 | https://www.pavolsurda.com/rhinosinusitis-in-children | Rhinosinusitis in Children | Pavol Surda
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Rhinosinusitis
in children and adolescents
What is paediatric rhinosinusitis?
You may hear and read different terms in relation to your child’s nasal symptoms. “Paediatric” refers to healthcare issues in childhood. “Rhinitis” is inflammation of the lining of the nose and can be caused in the short term by common infections, or over a longer period by an allergy. “Rhinosinusitis” is inflammation that affects not only the nose but also the sinuses. This can be chronic (lasting more than 12 weeks) or acute (less than 12 weeks).
What is the difference between acute and chronic rhinosinusitis?
Acute rhinosinusitis (ARS) is the sudden onset of two or more of the following symptoms;
Nasal blockage/obstruction/congestion
Or discoloured nasal discharge
Or cough (daytime and night-time)
Strictly speaking, acute rhinosinusitis includes episodes of these symptoms that last up to 12 weeks; in reality, most infections last only a few weeks and the majority settle down spontaneously.
Chronic rhinosinusitis (CRS) causes similar symptoms (listed below) but is present for 12 weeks or more:
Two or more of the below symptoms, one of which should be either
Nasal blockage / obstruction / congestion
or nasal discharge (anterior/posterior nasal drip)
+/- facial pain/pressure;
+/-cough.
It can be more challenging to diagnose CRS in children, partly because they find it difficult to describe and communicate their symptoms. In addition, it can be hard to tell the difference between CRS and other common reasons for a blocked nose in children (such as allergic conditions, and enlarged adenoids). Your doctor will therefore take more information and examine your child to try to understand which condition your child has.
Chronic rhinosinusitis (CRS)
How common is it?
It is difficult to say exactly how common chronic rhinosinusitis is in children but studies have shown between 2.1 to 4% of children may have symptoms associated with sinus disease. This makes it a less common problem than in adults, but there is undoubtedly a significant impact on the quality of life for children who suffer from CRS. The impact on overall health is greater than that of diseases such as asthma, attention deficit hyperactivity disorder, juvenile rheumatoid arthritis, and epilepsy.
What causes this?
The factors that lead to chronic rhinosinusitis in children are not fully understood. Overall, it seems most likely that rather than one single cause, it is due to a combination of several different factors that result in inflammation.
These factors are not equally present in all children, and your child may only have some of the factors described below:
Cigarette smoking (Passive smoking)
This is one thing that we do know contributes. One study showed that 68% of children with symptoms of acute rhinosinusitis were exposed to passive smoking, compared to 1.2% among children without exposure
Studies have also demonstrated worse outcomes in children with CRS exposed to cigarette smoke. These include the need for more operations.
Adenoids (a mound of immune tissue, similar to a tonsil that lies at the back of the nose) can contribute to CRS in children by harbouring bacteria and by causing blockage of the nasal airway.
Although more clear evidence is needed there appear to be linked to allergy, asthma and reflux disease.
Rarely there are some conditions that can affect the nose and sinuses as well as the lungs such as cystic fibrosis that will require further investigation and treatment.
What tests for CRS might a doctor offer for my child? Depending on the information your doctor has as well as the examination findings there may be some further tests required.
These may include:
Allergy testing (skin prick or blood tests)
Scans of the nose and sinuses
Assessment of the breathing/lungs
Genetic testing
What treatments for CRS might a doctor offer for my child? The treatment of chronic rhinosinusitis in children is largely based on the therapies that have been found to be effective in adults. There is no good evidence in the literature to support the use of antibiotics for chronic rhinosinusitis in children.
The main medical treatments are:
Nasal steroid medication
Saline irrigation
If nasal steroids are to be used, then your doctor will choose one which has as few side effects as possible and generally this is a very safe way to use medication. Often, they will ask to measure your child’s height and weight regularly to make sure that the medication is not causing any problems.
Surgery can also be used to improve the control of symptoms in children, but usually, this tends to be reserved for only the very few cases that do not respond to medications. A surgeon will usually choose the least invasive possible type of surgery for your child in order to balance the risks of surgery against improving the symptoms of the disease.
Nasal polyps
It is rare for children to develop nasal polyps.
Some children with allergic rhinitis may develop swollen turbinates (part of the internal lining of the nose) that may be mistaken for polyps. If your child does have nasal polyps these may just be caused by inflammation but it is likely that your doctor will request further tests to look for a cause. Cystic fibrosis, a relatively rare inherited disorder that affects the respiratory system is commonly associated with chronic rhinosinusitis with nasal polyps.
Acute rhinosinusitis
Acute rhinosinusitis (ARS) as defined at the beginning of this leaflet covers a range of conditions from the common cold to bacterial rhinosinusitis with associated complications. What is important to understand is how rare symptoms of a blocked, runny nose will turn into anything other than the common cold.
It is believed that school children suffer from around 7-10 episodes of the common cold every year. Similarly to CRS, there is little in the way of scientific evidence to tell us who is more likely to get ARS. We do know that smoking and passive smoking increase the chances of developing ARS.
There are many different viruses that can cause the common cold some of these include:
Rhinovirus
Respiratory syncytial virus
Influenza virus
Adenovirus
The diagram below shows how the severity of the symptoms and the length of time assist us in making the diagnosis.
The common cold lasts typically 5-10 days but after 48 hours your child should experience a gradual improvement in symptoms. Post viral ARS either has symptoms lasting longer than 10 days or an increase in symptoms after 5 days.
Bacterial rhinosinusitis is rare but it is defined by the measures set out in the diagram such as a high temperature, severe pain, one-sided symptoms, getting better than getting worse (double sickening) and abnormal blood tests.
Can ARS make my child very unwell?
Complications of ARS are uncommon but vital to identify and if you suspect any of these you should speak to a medical professional immediately.
They most often occur early in the course of the illness, and some of the signs and symptoms to watch out for are:
Swelling or redness around the eye
Double vision
Lethargy and confusion
Severe headache or swelling over the forehead
Rash, difficulty in bright lights or neck stiffness
Current evidence suggests that antibiotic treatment of ARS in general practice does not prevent complications.
Hopkins C, Surda P, Walker A, Wolf A, Speth M, Jacques T, et al. EPOS 4 patients. Rhinology. 2021 Suppl. 30: 1-57.
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13783 | https://math.stackexchange.com/questions/502571/show-that-gcda-b-operatornamelcma-b-if-and-only-if-a-b | elementary number theory - Show that $\gcd(a,b)=\operatorname{lcm}(a,b)$ if and only if $a=b$. - Mathematics Stack Exchange
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Show that \gcd(a,b)=\operatorname{lcm}(a,b) if and only if a=b.
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\begingroup
I know how to prove a=b only if \gcd(a,b)=\operatorname{lcm}(a,b), but I don't know how to prove the "if part". Can anyone help me?
elementary-number-theory
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edited Sep 29, 2013 at 17:23
user26857
asked Sep 23, 2013 at 16:17
MaxGaussianMaxGaussian
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\begingroup What are you allowed to assume? For example, if you are allowed to use the fact that every number has a unique prime factorization, then the proof is almost trivial (if it still isn't, I can write it as an answer)\endgroup zodiac –zodiac 2013-09-23 16:21:00 +00:00 Commented Sep 23, 2013 at 16:21
\begingroup Have you attempted it by assuming that \gcd(a,b)\ne lcm(a,b)?\endgroup abiessu –abiessu 2013-09-23 16:23:33 +00:00 Commented Sep 23, 2013 at 16:23
2 \begingroup Note that we need to assume in addition that a and b are positive.\endgroup André Nicolas –André Nicolas 2013-09-23 16:33:03 +00:00 Commented Sep 23, 2013 at 16:33
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Hint:\gcd(a,b)\le a,b and {\rm lcm} (a,b)\ge a,b.
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answered Sep 23, 2013 at 16:30
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Let a<b
Then: \gcd(a,b)≤a\operatorname{lcm}(a,b)≥b
Let \gcd(a,b)=\operatorname{lcm}(a,b)=x since \gcd(a,b)≤a≤b≤\operatorname{lcm}(a,b) so x≤a≤b≤x so a=b=x
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edited Sep 29, 2013 at 17:28
answered Sep 23, 2013 at 16:30
user93089user93089
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\begingroup Since we assumed a<b, then how can we write a \leq b in the \gcd(a,b) \leq a \leq b \leq lcm(a,b) ?\endgroup alu –alu 2020-08-22 23:32:29 +00:00 Commented Aug 22, 2020 at 23:32
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Hint: If (a,g)\in \mathbb{N}^{}\times \mathbb{N}^{}, a divides g and g divides a, then a=g.
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answered Sep 23, 2013 at 16:30
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According to the fundamental theorem of arithmetic, for any a,b we can write a=\prod_n p_n^{\alpha_n} ,b=\prod_n p_n^{\beta_n}, where {p_n } is the increasing sequence of primes. Now, gcd(a,b)=lcm(a,b) \Leftrightarrow \prod_n p_n^{\min(\alpha_n,\beta_n)}=\prod_n p_n^{\max(\alpha_n,\beta_n)} \Leftrightarrow \forall n \; \alpha_n=\beta_n \Leftrightarrow a=b.
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answered Sep 23, 2013 at 16:52
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A brain-dead approach to this problem would be to analyze the problem one prime at a time. Any positive integer can be uniquely represented by the number of times each prime divides that integer (i.e. its "prime factorization"). The \text{lcm} and \gcd are easy to express in this form....
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answered Sep 23, 2013 at 16:54
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\begingroup I'm sure this will work, but sometimes these results about gcd and lcm are used as lemmas to prove the fundamental theorem of arithmetic\endgroup zodiac –zodiac 2013-09-23 17:02:10 +00:00 Commented Sep 23, 2013 at 17:02
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0In the following answer, since we assumed a<b, then how can we write a \leq b in \gcd(a,b)\leq a\leq b\leq\operatorname{lcm}(a,b) statement?
Related
1Some Number Theory proofs about GCD and LCM
3Relating \operatorname{lcm} and \gcd using inverses or fractions
11\gcd(a,\operatorname{lcm}(b,c))=\operatorname{lcm}(\gcd(a,b),\gcd(a,c)),\, gcd distributes over lcm
1\operatorname{lcm}(\operatorname{gcd}(a,b),c)=\operatorname{gcd}(\operatorname{lcm}(a,c),\operatorname{lcm}(b,c)) for any a,b,c \in \mathbb{Z}
0Problem with equation containing \operatorname{lcm} and \gcd
5Prove that if a and b are positive integers satisfying \gcd(a,b)=\operatorname{lcm}(a,b),then a=b
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13784 | https://www.pdesas.org/ContentWeb/Content/Content/13382/Lesson | Scientific Notation - SAS
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Scientific Notation
Lesson
Scientific Notation
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Grade Levels 8th Grade
Related Academic StandardsCC.2.2.8.B.1 Apply concepts of radicals and integer exponents to generate equivalent expressions.
Assessment AnchorsM08.B-E.1 Demonstrate an understanding of expressions and equations with radicals and integer exponents.
Eligible ContentM08.B-E.1.1.4 Perform operations with numbers expressed in scientific notation, including problems where both decimal and scientific notation are used. Express answers in scientific notation and choose units of appropriate size for measurements of very large or very small quantities (e.g., use millimeters per year for seafloor spreading). Interpret scientific notation that has been generated by technology (e.g., interpret 4.7EE9 displayed on a calculator as 4.7 × 10^9). M08.B-E.1.1.3 Estimate very large or very small quantities by using numbers expressed in the form of a single digit times an integer power of 10 and express how many times larger or smaller one number is than another. Example: Estimate the population of the United States as 3 × 10^8 and the population of the world as 7 × 10^9 and determine that the world population is more than 20 times larger than the United States’ population.
Big Ideas
Concepts
Competencies
Objectives
Students will learn to use exponents appropriately. Students will:
represent very large and very small numbers using scientific notation, including the relationship between positive and negative exponents and place value.
use scientific notation as a tool in representing, calculating, and solving expressions with real-world data.
Essential Questions
How is mathematics used to quantify, compare, represent, and model numbers?
Vocabulary
Exponent: A numeral that tells how many times a number or variable is used as a factor. For example, in 2 7, 2 is the base and 7 is the exponent; this means 2 is multiplied by itself 7 times.
Scientific Notation: A way of writing a number of terms of an integer power of
10 multiplied by a number greater than or equal to 1 and less than 10.
Duration
90–120 minutes
Prerequisite Skills
Prerequisite Skills haven't been entered into the lesson plan.
Materials
note cards with numbers on them (1–9 and many with zeros)
Scientific Notation worksheet (M-8-4-2_Scientific Notation and KEY.doc) for each student
Activity Sheet (M-8-4-2_Activity Sheet and KEY.doc) for each student
Optional Quiz (M-8-4-2_Optional Short Quiz and KEY.docx) for any student who may benefit
Related Unit and Lesson Plans
Exponents and Radicals
Exponents and Their Properties
Square and Cube Roots
Related Materials & Resources
The possible inclusion of commercial websites below is not an implied endorsement of their products, which are not free, and are not required for this lesson plan.
Scientific Notation—Practice converting between scientific notation and normal numbers.
Exponents: Scientific Notation
Scientific Notation
Decimals to Scientific Notation
King Kong Scientific Notation Game
Formative Assessment
View
Scientific Notation worksheet requires students to translate between standard notation and scientific notation. Plan necessary re
Teaching according to the types of errors made by students.
Activity worksheet indicates how well students are able to use their knowledge of place value to represent parts per thousand in terms of scientific notation.
The Optional Quiz (M-8-4-2_Optional Short Quiz and KEY.docx) may also be used to assess student mastery.
Suggested Instructional Supports
View Active Engagement, Modeling, Explicit Instruction, Formative Assessment
W:This lesson engages students in scientific notation. Students will change numbers between scientific notation and standard notation, do computations, and make representations of scientific notation with calculators.
H:Hook students by presenting scientific notation as a useful tool that will make their mathematical work easier, quicker, more efficient, and less error prone. Having students perform long computational tasks with both standard and scientific notation gives them firsthand experience in finding the advantages of using it proficiently.
E:Students review place value and the relationship between each place value’s name and association with one exponential power of ten. In this way, students use something familiar to get to the next level in their understanding of how to use scientific notation. Students can count each iteration of multiplying by ten and relate it to the exponent of the expression.
R:Students review lesson concepts in Activity 3, where actual data and governmental standards are used in a real-world application of scientific notation. By relating air quality data to specific objects such as particles, students can associate it with something that has an impact on their quality of life.
E:By completing a scientific notation worksheet, students will demonstrate their competence in applying their understanding of the relationship between scientific notation and standard notation, use skills in converting from one to the other, and evaluate the reasonableness of their answers.
T:The Extension may be used to tailor the lesson to meet the needs of students. The Routine section provides suggestions for reviewing lesson concepts throughout the year. The Small Group section is intended for students who may benefit from additional practice. The Expansion section includes ideas for students who are prepared for a challenge beyond the requirements of the standard.
O:This lesson begins with a discussion about very small and very large numbers. It introduces scientific notation as a way to represent these numbers beginning with a review of place value. Place value is then related to powers of ten. Scientific notation is introduced, and students are taught a few methods to determine how to represent large and small numbers with the looping method and calculator notation. A real-world example of air quality reinforces this information.
Instructional Procedures
View
Write a number on the chalkboard that contains at least ten digits. Ask students to multiply that number by 2, 3, and 10. After a few minutes ask whether students found this task difficult or time consuming. “There is a way to represent very small and very large numbers to complete calculations more quickly. Scientific notation allows us to write these numbers in a shorter, more manageable form.”
This section will review place value and powers of ten. Use the number written on the board and have volunteers come to the board and write the place value of any number.
Example: 1,324,890,625
1 – billions
3 – hundred millions
2 – ten millions
4 – millions
8 – hundred thousands
9 – ten thousands
0 – thousands
6 – hundreds
2 – tens
5 – ones
“At your desk, please write how many zeros each place value has following it. In other words, if there were only billions, millions, . . . how many zeros would there be?” Allow 5–10 minutes for students to work. If some finish early, ask them to help students who might need additional practice.
1 – billions (9 zeros)
,
3 – hundred millions (8 zeros)
2 – ten millions (7 zeros)
4 – millions (6 zeros)
,
8 – hundred thousands (5 zeros)
9 – ten thousands (4 zeros)
0 – thousands (3 zeros)
,
6 – hundreds (2 zeros)
2 – tens (1 zero)
5 – ones (0 zeros)
“These place values all relate to powers of ten. Take a few minutes to do the following multiplication problems.” Distribute short worksheets with multiplication problems that have 10 multiplied by 10 numerous times or write 15 problems on the board.
Example:
10 ×10 = 10 2 = 100
10 ×10 ×10 = 10 3 = 1,000
10 ×10 ×10 ×10 = 10 4 = 10,000
Students can calculate by hand or use calculators. If they use calculators, have them copy their result onto their paper.
“Did you find any patterns with the number of zeros each answer had and the number of times you multiplied 10?”Students should respond that the number of times they multiplied by ten was the number of zeros. “The following chart shows how place value relates to powers of 10.”
Project the chart:
Place Value and Powers of 10
Place ValueStandard FormNumber of ZerosExponential Form
…………
Hundred Billions 100,000,000,000 11 10 11
Ten Billions 10,000,000,000 10 10 10
Billions 1,000,000,000 9 10 9
Hundred Millions 100,000,000 8 10 8
Ten Millions 10,000,000 7 10 7
Millions 1,000,000 6 10 6
Hundred Thousands 100,000 5 10 5
Ten Thousands 10,000 4 10 4
Thousands 1,000 3 10 3
Hundreds 100 2 10 2
Tens 10 1 10 1
Ones 1 0 10 0
Tenths 0.1 1 10−1
Hundredths 0.01 2 10−2
Thousandths 0.001 3 10−3
Ten Thousandths 0.0001 4 10−4
Hundred Thousandths 0.00001 5 10−5
…………
“Notice that for each place value, the number of zeros it has in standard form becomes the power of 10 when the place value is written in exponential form. For values above 1, the exponent is positive. But for values below 1 (our decimal place values), the exponent is negative.
“As you can see, any place value can be expressed as a power of 10. This means that if we are dealing with a very large or very small number (or in other words, a number with many digits), we may abbreviate the number by expressing it as the product of an integer and a power of ten. This way of expressing a number is called scientific notation.
“A number is expressed in scientific notation when it is written as the product of a factor and a power of 10. The factor must be greater than or equal to 1 and less than 10.”
Write on the board:
Scientific Notation
a × 10 n, where 1 ≤ a< 10 and n is an integer
“Let’s look at an example. We’ll start with a large number, like 5 million.
“What is the largest place value in 5 million?”(millions) “How many zeros are in a million?”(6) “How can we express millions as a power of 10?”(10 6) “So if 10 6 represents 1 million, what should we multiply 1 million by to get 5 million?”(5)
5 million = 5,000,000 = 5 1,000,000 =
Activity 1
“Please write the following number in scientific notation. As you try to do this, consider the numbers of zeros that follow each number and which power of 10 that represents.”
7,000,000,000 (7 ×10 9)
50,000 (5 × 10 4)
20 (2 × 10 1)
0.0006 (6 ×10−4)
0.8 (8 × 10−1)
900,000,000,000,000 (9 ×10 14)
“Despite our recent examples, there are other numbers between 1 and 10 that are not whole numbers. We can also use decimals as the factors in scientific notation. For example, let’s consider the number 5,930,000,000. What is the highest place value in this number?”(billions)“According to our place value chart and what we have learned about corresponding powers of 10, what power of 10 is associated with one billion?”(9)“This implies that our scientific notation will look like: __ × 10 9. So the only issue now is to identify the factor. This time, we cannot simply say that the factor is 5, because 5 × 10 9 = 5 billion (or 5,000,000,000). Yet we want 5,930,000,000. If we call the factor 6, then we have 6 × 10 9 = 6 billion (or 6,000,000,000), and this is too big. So our factor must be between 5 and 6. Is it 5.1, 5.2, 5.3, . . . ?”Guide students toward realizing that the factor must be 5.93.
5,930,000,000 = 5.93 × 10 9
“Please convert the following numbers from standard notation to scientific notation.”
617 (6.17 _×_ 10 2)
9,125,600,000,000,000 (9.1256 _×_ 10 15)
0.000345 (3.45 × 10−4)
[Note that the LARGEST place value in 0.000345 is the ten thousandths, NOT the thousand thousandths, as 0.0001 > 0.0000001.]
The Looping Method
This system will help students understand how they are moving the decimal place. “Now we will look at a slightly different strategy when it comes to converting between standard and scientific notation. In this strategy, we will use the exponent to count how many times to move the decimal place.” Write the following example on the board:
5,032,000
“First, locate the position of the decimal point in the standard form of the number.” (5,032,000.)“Now, we will move the decimal point until the number is between 1 and 10. What direction will we need to move the decimal point in order for 5,032,000 to become a number in between 1 and 10?”(left) “Start moving the decimal point left. Keep going until the number you see is between 1 and 10.”
Think: Is 5,032,000 between 1 and 10? If no, keep going.
Think: Is 503,200 between 1 and 10? If no, keep going.
Think: Is 50,320 between 1 and 10? If no, keep going.
Think: Is 5,032 between 1 and 10? If no, keep going.
Think: Is 503.2 between 1 and 10? If no, keep going.
Think: Is 50.32 between 1 and 10? If no, keep going.
Think: Is 5.032 between 1 and 10? If yes, stop!
“By moving the decimal point left 1 place at a time, we eventually get to the number 5.032. This value is indeed greater than or equal to 1 and less than 10. This means that 5.032 is the factor in our scientific notation: 5.032 × 10?. Now all we have to do is determine the exponent. To do so, think about how many times we needed to move the decimal point to get from 5,032,000 to 5.032. How many times was this?”(6 to the left)“Since we moved the decimal point 6 times to the left, our exponent is +6.”
5,032,000 = 5.032 × 10 6
“Let’s look at another example.”
0.002705
“First, locate the position of the decimal point when the number is written in standard form.” (0.002705) “Now, we will move the decimal point until the number we see is between 1 and 10. What direction will we have to move the decimal point to accomplish this?”(right) “Start moving the decimal point to the right one place value at a time. Don’t stop until you see a number between 1 and 10.”
Think: Is 0.002705 between 1 and 10? If no, keep going.
Think: Is 0.02705 between 1 and 10? If no, keep going.
Think: Is 0.2705 between 1 and 10? If no, keep going.
Think: Is 2.705 between 1 and 10? If yes, stop!
“By moving the decimal point right one place at a time, we eventually get to the number 2.705. This value is indeed greater than or equal to 1 and less than 10. This means that 2.705 is the factor in our scientific notation: 2.705 × 10?. Now all we have to do is determine the exponent. To do so, think about how many times we needed to move the decimal point to get from 0.002705 to 2.705. How many times was this?”(3 to the right)“Since we moved the decimal point 3 times to the right, our exponent is −3.”
0.002705 = 2.705 × 10−3
“The looping method can also be helpful when starting with a number in scientific notation and converting it to standard form. All we have to do is go backwards.”
3.8 × 10 7
“Once again, we start by locating the position of the decimal point.” (3.8 × 10 7) “This time, however, instead of taking a very large number and moving the decimal point so it becomes a number in between 1 and 10, we are starting with a factor in between 1 and 10 and moving the decimal point back so it is a very large number. What direction will we need to move the decimal point to create a very large number?” (right) “How many times will we have to move the decimal point to the right?”(7, because this is the value of the exponent)
3.8 × 10 7 = 38,000,000
“Let’s convert another number from scientific notation into standard form.”
1.52 × 10−5
“As always, start by locating the position of the decimal point.” (1 .52 × 10−5) “The negative exponent indicates that this is a very small number. Which direction will we have to move the decimal point in 1.52 to convert this back to a very small number?”(left) “How many times will we have to move the decimal point to the left?” (5, because this is the value of the exponent)
1.52 × 10−5 = 0.0000152
Activity 2
Divide the class into two groups. Each student in one group will get a piece of paper with a zero printed on it. Students in the second group will have a piece of paper with one number, 1 through 9. Show the class a number written in scientific notation. Have the group with whole numbers decide who to designate as the beginning digit of the standard form of the number. Have the second group determine how many of their members it will take to complete the standard form of the number. Have these students stand in a row to display the result.
Alternative ways to play:
Use negative exponents and have the factor group be in charge of the decimal point.
Give students just the factor and the exponent without the full expression.
As the class progresses, have the factor include tenths and hundredths, reinforcing that the factor must be between 0 and 10.
Calculator Notation
Note: Calculator display notation varies widely among different brands of calculators and different purposes and capacities of the calculator. Bring some examples of as many different ones as you can find to the classroom to show students directly. Scientific notation has a very specific representation on calculators that is different than handwritten or printed text. Have students choose a whole number 2 through 9 and have them multiply that number repeatedly on their calculators. Have students raise their hands when more than just numbers appear on the calculator.
“Calculators do not show scientific notation in the form that we have been learning. Look at your calculators. Does anyone have a number like this?”
5.477E8
“Thinking about what we have already learned, how can we write this number in scientific notation?”
5.477 is the factor
8 is the exponent
5.477 ×10 8 = 547,700,000
“Other calculators have the factor in the main screen with the exponent of ten in the upper right.”
1.38412872 10
1.38412872 is the factor
10 is the power of ten
1.38412872 ×10 10 = 13,841,287,200
Activity 3
This activity will reinforce converting numbers into scientific notation in a real-world setting. Information is from the Pennsylvania Department of Environmental Protection (
“The Commonwealth of Pennsylvania has adopted air quality standards that regulate the amount of harmful pollutants in the air we breathe. I will pass the following table to each student so that we may study the standards together.” Hand out the Activity Sheet (M-8-4-2_Activity Sheet and KEY.doc). “Reviewing the standards, each pollutant is allowed in ppm, or parts per million. This means that for every million particles of air, the number listed is the amount of each pollutant allowed in those one million particles. We are going to convert the numbers to scientific notation, imagining what part of one particle of air can be made up of each pollutant.” Write the equations on the board as the answers are calculated.
“The standard for carbon monoxide is 9 ppm.
“Nine particles of carbon monoxide per million particles of air + carbon monoxide is
9 parts per million”
= 0.000001 or one millionth
9 ×0.000001 = 0.000009
0.000009 = 9 ×10−6
“So for each particle of air + carbon monoxide, 9 ×10−6 particles may be carbon monoxide.
The standard for nitrogen dioxide is 0.053 ppm or 0.053 particles per million.”
0.053 = 5.3 ×10−2
= 0.000001 or one millionth
0.000001 = 1 ×10−6
“So we are looking for 0.053 0.000001.”
0.053 (1 ×10−6) = ?
“We have to make sure our factor is between 1 and 10. So if we move the decimal two places to the right, we have to make the exponent two smaller.
“So for each particle of air, 5.3 ×10−8 of it can be nitrogen dioxide.
“Complete the chart. If you have difficulties, consult your neighbor for help or raise your hand.
“Remember the rule for multiplying the same bases with exponents: Add the exponents. For example, 10 3 × 10−4 = 10−1 because 10(3+−4) = 10(−1).”
Use the Scientific Notation worksheet (M-8-4-2_Scientific Notation and KEY.doc). To reinforce and evaluate the concepts introduced, list equations that look like scientific notation, but are not written correctly. Have students find the correct equations and explain why the original numbers were written incorrectly.
Examples:
0.002 ×10−6
5326 ×10−2
984 ×10 9
64.1 ×10 0
0.3 ×10 6
53 ×10 1
Solutions:
2 ×10−9
5.326 ×10 1
9.84 ×10 11
6.41 ×10 1
3 ×10 5
5.3 ×10 2
Allow students to work alone or in pairs to find the correct expressions. When most are finished, have students write their work on the board. Go over the problems as a class. Students should respond that each of these is incorrect because the factor is not between 1 and 10, and therefore must be adjusted to be properly written in scientific notation.
Extension:
Routine:Have students find their own real−world examples of the use of very small and very large numbers. Very large numbers can include the population of countries, the number of fish in the ocean, or even the amount of money in the stock exchange. Very small numbers might include the size of atoms or water quality requirements. Encourage students to find their own examples and write these numbers in scientific notation. Look for the opportunity to include scientific notation in journal writing in other subject areas.
Students may also review conversion using one of the games at the following web addresses:
Small Group:Students who may benefit from additional practice with scientific notation and decimal conversion may be allowed to play the online King Kong game at the
websites listed below:
If students require additional instruction, the following websites may be used:
Expansion: Students who are prepared for a challenge beyond the requirements of the standard may be given the activities below.
Have students write the following expressions in scientific notation.
5.6 ×10 6 ×10−9
2.39 ×10 2 ×10 5
9× 10−3 ×10−1
87.75 ×10−6 ×10−3
Solutions for Expansion:
5.6 ×10−3
2.39 ×10 7
9 ×10−4
8.775 ×10−8
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9.15
Solving Rational Equations using the LCD
Written by:Lori Jordan | Kate Dirga
Fact-checked by:The CK-12 Editorial Team
Last Modified: Sep 01, 2025
A right triangle has leg lengths of 12 and 1x units. Its hypotenuse is 2 units. What is the value of x?
Solving Rational Equations
In addition to using cross-multiplication to solve a rational equation, we can also use the LCD of all the rational expressions within the equation and eliminate the fraction. To demonstrate, we will walk through a few problems.
Let's solve the following rational equations.
52+1x=3
The LCD for 2 and x is 2x. Multiply each term by 2x, so that the denominators are eliminated. We will put the 2x over 1, when multiplying it by the fractions, so that it is easier to line up and cross-cancel.
52+1x=32x1⋅52+2x1⋅1x=2x⋅35x+2=6x2=x
Checking the answer, we have 52+12=3→62=3
5xx−2=7+10x−2
Because the denominators are the same, we need to multiply all three terms by x−2.
5xx−2=7+10x−2(x−2)⋅5xx−2=(x−2)⋅7+(x−2)⋅10x−25x=7x−14+10−2x=−4x=2
Checking our answer, we have: 5⋅22−2=7+102−2→100=7+100. Because the solution is the vertical asymptote of two of the expressions, x=2 is an extraneous solution. Therefore, there is no solution to this problem.
3x+45=6x−2
Determine the LCD for 5, x, and x−2. It would be the three numbers multiplied together: 5x(x−2). Multiply each term by the LCD.
3x+45=6x−25x(x−2)1⋅3x+5x(x−2)1⋅45=5x(x−2)1⋅6x−215(x−2)+4x(x−2)=30x
Multiplying each term by the entire LCD cancels out each denominator, so that we have an equation that we have learned how to solve in previous concepts. Distribute the 15 and 4x, combine like terms and solve.
15x−30+4x2−8x=30x4x2−23x−30=0
This polynomial is not factorable. Let’s use the Quadratic Formula to find the solutions.
x=23±√(−23)2−4⋅4⋅(−30)2⋅4=23±√10098
Approximately, the solutions are 23+√10098≈6.85 and 23−√10098≈−1.096. It is harder to check these solutions. The easiest thing to do is to graph 3x+45 in Y1 and 6x−2 in Y2 (using your graphing calculator).
The x-values of the points of intersection (purple points in the graph) are approximately the same as the solutions we found.
Examples
Example 1
Earlier, you were asked to find the value of x.
We need to use the Pythagorean Theorem to solve for x.
(12)2+(1x)2=2214+1x2=44x21⋅14+4x21⋅1x2=4⋅4x2x2+4=16x24=15x2415=x2x=2√1515
Solve the following equations.
Example 2
2xx−3=2+3xx2−9
The LCD is x2−9. Multiply each term by its factored form to cross-cancel.
2xx−3=2+3xx2−9(x−3)(x+3)1⋅2xx−3=(x−3)(x+3)⋅2+(x−3)(x+3)1⋅3xx2−92x(x+3)=2(x2−9)+3x2x2+6x=2x2−18+3x3x=−18x=−6
Checking our answer, we have: 2(−6)−6−3=2+3(−6)(−6)2−9→−12−9=2+−1827→43=2−23
Example 3
4x−3+5=9x+2
The LCD is (x−3)(x+2). Multiply each term by the LCD.
4x−3+5=9x+2(x−3)(x+2)⋅4x−3+(x−3)(x+2)⋅5=(x−3)(x+2)⋅9x+24(x+2)+5(x−3)(x+2)=9(x−3)4x+8+5x2−5x−30=9x−275x2−10x+5=05(x2−2x+1)=0
This polynomial factors to be 5(x−1)(x−1)=0, so x=1 is a repeated solution. Checking our answer, we have 41−3+5=91+2→−2+5=3
Example 4
3x2+4x+4+1x+2=2x2−4
The LCD is (x+2)(x+2)(x−2).
3x2+4x+4+1x+2=2x2−4(x+2)(x+2)(x−2)⋅3(x+2)(x+2)+(x+2)(x+2)(x−2)⋅1x+2=(x+2)(x+2)(x−2)⋅2(x−2)(x+2)3(x−2)+(x−2)(x+2)=2(x+2)3x−6+x2−4=2x+4x2+x−14=0
This quadratic is not factorable, so we need to use the Quadratic Formula to solve for x.
x=−1±√1−4(−14)2=−1±√572≈3.27 and −4.27
Using your graphing calculator, you can check the answer. The x-values of points of intersection of y=3x2+4x+4+1x+2 and y=2x2−4 are the same as the values above.
Review
Determine if the following values for x are solutions for the given equations.
4x−3+2=3x+4, x=−1
2x−1x−5−3=x+62x, x=6
What is the LCD for each set of numbers?
4−x, x2−16
2x, 6x−12, x2−9
x−3, x2−x−6, x2−4
Solve the following equations.
6x+2+1=5x
53x−2x+1=4x
12x2−9=8xx−3−2x+3
6xx2−1+2x+1=3xx−1
5x−34x−x+1x+2=1x2+2x
4xx2+6x+9−2x+3=3x2−9
x2x2−8x+16=xx−4+3xx2−16
5x2x−3+x+1x=6x2+x+122x2−3x
3xx2+2x−8=x+1x2+4x+2x+1x2−2x
x+1x2+7x+x+2x2−3x=xx2+4x−21
Review (Answers)
Click HERE to see the answer key or go to the Table of Contents and click on the Answer Key under the 'Other Versions' option.
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13786 | https://analystprep.com/cfa-level-1-exam/quantitative-methods/properties-of-continous-uniform-distribution/ | Save 10% on All AnalystPrep 2024 Study Packages with Coupon Code BLOG10.
Properties of Continuous Uniform Distribution
The continuous uniform distribution is such that the random variable (X) takes values between (a) (lower limit) and (b) (upper limit). In the field of statistics, (a) and (b) are known as the parameters of continuous uniform distribution. We cannot have an outcome of either less than (a) or greater than (b).
The probability density function for this type of distribution is:
$$ { f }_{ x }\left( x \right) =\frac { 1 }{ b -a } \quad \quad a< x < b $$
(X \sim U (a,b)) is the most commonly used shorthand notation read as “the random variable (x). It has a continuous uniform distribution with parameters (α) and (β).”
The total probability is spread uniformly between the two limits. Intervals of the same length have the same probability.
In other words, for all (a \le x_1 < x_2 \le b), we have
$$P(X < a) = P(X > b = 0)$$
And the probability that the random variable will have a value between (x_1) and (x_2) is given as follows:
$$P(x_1 \le X \le x_2) =\cfrac {(x_2 – x_1)}{(b – a)}$$
The mean and variance of continuous uniform distribution is given below:
$$ \text{Mean} =\cfrac {(a + b)}{2} $$
$$ \text{Variance} =\cfrac {(b – a)^2}{12} $$
Example: Probability Density Function
You have been given that (Y \sim U(100,300)).
Calculate (P(Y > 174)) and (P(100 < Y < 226).
Solution
The probability density function is given by:
$$ f_x(x) =\cfrac {1}{(300 – 100)} =\cfrac {1}{200} $$
Therefore, each “unit interval” has a probability of (\frac {1}{200}).
This means that (P(Y > 174) =\cfrac {(300 – 174)}{200} = \cfrac {126}{200} = 0.63).
Similarly, (P(100 < Y < 226) = 0.63) because the interval has the same length as above (126) hence the same probability.
Cumulative Distribution Function of a Continuous Uniform Distribution
Remember that a cumulative distribution function, (F(x)), gives the probability that the random variable (X) is less than or equal to (x), for every (x) value. It is usually expressed as:
$$F(x)=P(X\leq x)$$
The cumulative distribution function of the continuous uniform distribution increases linearly from (α) to (β).
The CDF is linear over the variable’s range, and it is given by:
$$
F(x)=P(X\leq x)=\frac{x-a}{b-a}
$$
Question
A random variable (X) is uniformly distributed between 32 and 42. What is the probability that (X) will be between 32 and 40?
Solution
The correct answer is C.
First, you should determine the pdf:
$$ \begin{align} f_x(x) & =\cfrac {1}{(42 – 32)} \ &=\cfrac {1}{10} \ \end{align} $$
Therefore,
$$ P(32 < Y < 40) =\cfrac {(40 – 32)}{10} = 0.8 \text{ or } 80\% $$
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13787 | https://econ.tau.ac.il/sites/economy.tau.ac.il/files/media_server/Economics/PDF/seminars%202016-17/OSP.pdf | Obviously Strategy-Proof Mechanisms∗ Shengwu Li† Job Market Paper First uploaded: 3 February 2015. This version: 16 March 2016.
Abstract What makes some strategy-proof mechanisms easier to understand than others? To address this question, I propose a new solution concept: A mechanism is obviously strategy-proof (OSP) if it has an equilibrium in obviously dominant strategies. This has a behavioral interpretation: A strategy is obviously dominant if and only if a cognitively limited agent can recognize it as weakly dominant. It also has a classical interpretation: A choice rule is OSP-implementable if and only if it can be carried out by a social planner under a particular regime of partial commitment. I fully characterize the set of OSP mechanisms in a canonical setting, with one-dimensional types and transfers. A laboratory experiment tests and corroborates the theory.
1 Introduction Dominant-strategy mechanisms are often said to be desirable. They reduce participation costs and cognitive costs, by making it easy for agents to decide what to do.1 They protect agents from strategic errors.2 Dominant-strategy mechanisms prevent waste from rent-seeking espionage, since spying on other players yields no strategic advantage. Moreover, the resulting outcome does not depend sensitively on each agent’s higher-order beliefs.3 These benefits largely depend on agents understanding that the mechanism has an equilibrium in dominant strategies; i.e. that it is strategy-proof (SP). Only then can they conclude that they need not attempt to discover their opponents’ strategies or to game the system.4 ∗I thank especially my advisors, Paul Milgrom and Muriel Niederle. I thank Nick Arnosti, Douglas Bernheim, Gabriel Carroll, Paul J. Healy, Matthew Jackson, Fuhito Kojima, Roger Myerson, Michael Ostrovsky, Matthew Rabin, Alvin Roth, and Ilya Segal for their invaluable advice. I thank Paul J. Healy for his generosity in allowing my use of the Ohio State University Experimental Economics Laboratory. I thank Muriel Niederle and the Stanford Economics Department for financial support for the experiment. This work was supported by the Kohlhagen Fellowship Fund, through a grant to the Stanford Institute for Economic Policy Research. All errors remain my own.
†shengwu@stanford.edu 1Vickrey (1961) writes that, in second-price auctions: “Each bidder can confine his efforts and attention to an appraisal of the value the article would have in his own hands, at a considerable saving in mental strain and possibly in out-of-pocket expense.” 2For instance, school choice mechanisms that lack dominant strategies may harm parents who do not strategize well (Pathak and S¨ onmez, 2008).
3Wilson (1987) writes, “Game Theory has a great advantage in explicitly analyzing the consequences of trading rules that presumably are really common knowledge; it is deficient to the extent it assumes other features to be common knowledge, such as one player’s probability assessment about another’s preferences or information.” 4Policymakers could announce that a mechanism is strategy-proof, but that may not be enough. If agents do not understand the mechanism well, then they may be justifiably skeptical of such declarations. For instance, Google’s 1 However, some strategy-proof mechanisms are simpler for real people to understand than others.
For instance, choosing when to quit in an ascending clock auction is the same as choosing a bid in a second-price sealed-bid auction (Vickrey, 1961). The two formats are strategically equivalent; they have the same reduced normal form.5 Nonetheless, laboratory subjects are substantially more likely to play the dominant strategy under a clock auction than under sealed bids (Kagel et al., 1987). Theorists have also expressed this intuition: Some other possible advantages of dynamic auctions over static auctions are difficult to model explicitly within standard economics or game-theory frameworks. For example, . . . it is generally held that the English auction is simpler for real-world bidders to understand than the sealed-bid second-price auction, leading the English auction to perform more closely to theory. (Ausubel, 2004) In this paper, I model explicitly what it means for a mechanism to be obviously strategy-proof.
This approach invokes no new primitives. Thus, it identifies a set of mechanisms as simple to understand, while remaining as parsimonious as standard game theory.
A strategy Si is obviously dominant if, for any deviating strategy S′ i, starting from any earliest information set where Si and S′ i diverge, the best possible outcome from S′ i is no better than the worst possible outcome from Si. A mechanism is obviously strategy-proof (OSP) if it has an equilibrium in obviously dominant strategies. By construction, OSP depends on the extensive game form, so two games with the same normal form may differ on this criterion. Obvious dominance implies weak dominance, so OSP implies SP.
This definition distinguishes ascending auctions and second-price sealed-bid auctions. Ascending auctions are obviously strategy-proof. Suppose you value the object at $10. If the current price is below $10, then the best possible outcome from quitting now is no better than the worst possible outcome from staying in the auction (and quitting at $10). If the price is above $10, then the best possible outcome from staying in the auction is no better than the worst possible outcome from quitting now.
Second-price sealed-bid auctions are strategy-proof, but not obviously strategy-proof. Consider the strategies “bid $10” and “bid $11”. The earliest information set where these diverge is the point where you submit your bid. If you bid $11, you might win the object at some price strictly below $10. If you bid $10, you might not win the object. The best possible outcome from deviating is better than the worst possible outcome from truth-telling. This captures an intuition expressed by experimental economists: The idea that bidding modestly in excess of x only increases the chance of win-ning the auction when you don’t want to win is far from obvious from the sealed bid procedure. (Kagel et al., 1987) I produce two characterization theorems, which suggest two interpretations of obvious strategy-proofness. The first interpretation is behavioral: Obviously dominant strategies are those that can be recognized as dominant by a cognitively limited agent. The second interpretation is classical: Obviously strategy-proof mechanisms are those that can be carried out by a social planner with only partial commitment power.
advertising materials for the Generalized Second-Price auction appeared to imply that it was strategy-proof, when in fact it was not (Edelman et al., 2007). Moreover, Rees-Jones (2015) and Hassidim et al. (2015) find evidence of strategic mistakes in approximately strategy-proof matching markets, even though participants face a high-stakes decision with expert advice.
5This equivalence assumes that we restrict attention to cut-offstrategies in ascending auctions.
2 Figure 1: ‘Similar’ mechanisms from 1’s perspective.
First, I model an agent who has a simplified mental representation of the world: Instead of understanding every detail of every game, his understanding is limited by a coarse partition on the space of all games. I show that a strategy Si is obviously dominant if and only if such an agent can recognize Si as weakly dominant.
Consider the mechanisms in Figure 1. Suppose Agent 1 has preferences: A ≻B ≻C ≻D. In mechanism (i), it is a weakly dominant strategy for 1 to play L. Both mechanisms are intuitively similar, but it is not a weakly dominant strategy for Agent 1 to play L in mechanism (ii).
In order for Agent 1 to recognize that it is weakly dominant to play L in mechanism (i), he must use contingent reasoning. That is, he must think through hypothetical scenarios: “If Agent 2 plays l, then I should play L, since I prefer A to B. If Agent 2 plays r, then I should play L, since I prefer C to D. Therefore, I should play L, no matter what Agent 2 plays.” Notice that the quoted inferences are valid in (i), but not valid in (ii).
Suppose Agent 1 is unable to engage in contingent reasoning. That is, he knows that playing L might lead to A or C, and playing R might lead to B or D. However, he does not understand how, state-by-state, the outcomes after playing L are related to the outcomes after playing R. Then it is as though he cannot distinguish (i) and (ii).
This idea can be made formal and general. I define an equivalence relation on the space of mechanisms: The experience of agent i at history h records the information sets where i was called to play, and the actions that i took, in chronological order.6 Two mechanisms G and G′ are i-indistinguishable if there is a bijection from i’s information sets and actions in G, onto i’s information sets and actions in G′, such that: 1. G can produce for i some experience if and only if G′ can produce for i the image (under the bijection) of that experience.
2. An experience might result in some outcome in G if and only if its image might result in that same outcome in G′.
With this relation, we can partition the set of all mechanisms into equivalence classes. For instance, the mechanisms in Figure 1 are 1-indistinguishable.
6An experience is a standard concept in the theory of extensive games; experiences are sometimes used to define perfect recall.
3 Our agent knows the experiences that a mechanism might generate, and the resulting outcomes.
That is, he knows all the points at which he may be called to play, and all the actions available at each point. He knows, for any sequence of points he was called to play and actions that he took, whether the game might end and, if so, what outcomes might result. However, this does not nail down every detail of the mechanism. In particular, our agent is unable to reason contingently about hypothetical scenarios.
The first characterization theorem states: A strategy Si is obviously dominant in G if and only if it is weakly dominant in every G′ that is i-indistinguishable from G.
This shows that obviously dominant strategies are those that can be recognized as weakly dominant without contingent reasoning. An obviously dominant strategy is weakly dominant in any i-indistinguishable mechanism. In that sense, such a strategy is robustly dominant.
The second characterization theorem for OSP relates to the problem of mechanism design under partial commitment. In mechanism design, we usually assume that the Planner can commit to every detail of a mechanism, including the events that an individual agent does not directly observe. For instance, in a sealed-bid auction, we assume that the Planner can commit to the function from all bid profiles to allocations and payments, even though each agent only directly observes his own bid. Sometimes this assumption is too strong. If agents cannot individually verify the details of a mechanism, the Planner may be unable to commit to it.
Mechanism design under partial commitment is a pressing problem. Auctions run by central brokers over the Internet account for billions of dollars of economic activity (Edelman et al., 2007).
In such settings, bidders may be unable to verify that the other bidders exist, let alone what actions they have taken. As another example, some wireless spectrum auctions use computationally demanding techniques to solve complex assignment problems. In these settings, individual bidders may find it difficult and costly to verify the output of the auctioneer’s algorithm (Milgrom and Segal, 2015).
For the second characterization theorem, I consider a ‘metagame’ where the Planner privately communicates with agents, and eventually decides on an outcome. The Planner chooses one agent, and sends a private message, along with a set of acceptable replies. That agent chooses a reply, which the Planner observes.
The Planner can then either repeat this process (possibly with a different agent) or announce an outcome and end the game.
The Planner has partial commitment power: For each agent, she can commit to use only a subset of her available strategies. However, the subset she promises to Agent i must be measurable with respect to i’s observations in the game. That is, if the Planner plays a strategy not in that subset, then there exists some agent strategy profile such that Agent i detects (with certainty) that the Planner has deviated. We call this a bilateral commitment.
Suppose we require that each agent i’s strategy be optimal, for any strategies of the other agents, and for any Planner strategy compatible with the commitment made to i. What choice rules can be implemented in this metagame?
The second characterization theorem states: A choice rule can be supported by bilateral commit-ments if and only if that choice rule is OSP-implementable. Consequently, in addition to formalizing a notion of cognitive simplicity, OSP also captures the set of choice rules that can be carried out with only bilateral commitments.
After defining and characterizing OSP, I apply this concept to several mechanism design envi-ronments.
For the first application, I consider binary allocation problems. In this environment, there is a set of agents N with continuous single-dimensional types θi ∈[θi, θi]. An allocation y is a subset of N. An allocation rule fy is a function from type profiles to allocations. We augment this with a transfer rule ft, which specifies money transfers for each agent. Each agent has utility equal to 4 his type if he is in the allocation, plus his net transfer.
ui(θi, y, t) = 1i∈yθi + ti (1) Binary allocation problems encompass several canonical settings. They include private-value auctions with unit demand. They include procurement auctions with unit supply; not being in the allocation is ‘winning the contract’, and the bidder’s type is his cost of provision. They also include binary public good problems; the feasible allocations are N and the empty set.
Mechanism design theory has extensively investigated SP-implementation in this environment.
fy is SP-implementable if and only if fy is monotone in each agent’s type (Spence, 1974; Mirrlees, 1971; Myerson, 1981). If fy is SP-implementable, then the required transfer rule ft is essentially unique (Green and Laffont, 1977; Holmstr¨ om, 1979).
What are analogues of these canonical results, if we require OSP-implementation rather than SP-implementation? Are ascending clock auctions special, or are there other OSP mechanisms in this environment?
I prove the following theorem: Every mechanism that OSP-implements an allocation rule is ‘essentially’ a monotone price mechanism, which is a new generalization of ascending clock auctions.
Moreover, this is a full characterization of OSP mechanisms: For any monotone price mechanism, there exists some allocation rule that it OSP-implements.
These results imply that when we desire OSP-implementation in a binary allocation problem, we need not search the space of all extensive game forms. Without loss of generality, we can focus our attention on the class of monotone price mechanisms.7 Additionally, I characterize the set of OSP-implementable allocation rules.
For this part, I assume that the lowest type of each agent is never in the allocation, and is required to have a zero transfer.
I provide a necessary and sufficient condition for an allocation rule to be OSP-implementable.
As a second application, I produce an impossibility result for a classic matching algorithm: With 3 or more agents, there does not exist a mechanism that OSP-implements Top Trading Cycles (Shapley and Scarf, 1974).
I conduct a laboratory experiment to test the theory.
In the experiment, I compare three pairs of mechanisms. In each pair, both mechanisms implement the same allocation rule. One mechanism is obviously strategy-proof. The other mechanism is strategy-proof, but not obviously strategy-proof. Standard theory predicts that both mechanisms result in dominant strategy play, and have identical outcomes. Instead, subjects play the dominant strategy at significantly higher rates under the OSP mechanism, compared to the mechanism that is just SP. This effect occurs for all three pairs of mechanisms, and persists even after playing each mechanism five times with feedback.
The rest of the paper proceeds in the usual order. Section 2 reviews the literature. Section 3 provides formal definitions and characterizations. Section 4 covers applications. Section 5 reports the laboratory experiment. Section 6 concludes. Proofs omitted from the main text are in Appendix A.
2 Related Literature It is widely acknowledged that ascending auctions are simpler for real bidders than second-price sealed-bid auctions (Ausubel, 2004). Laboratory experiments have investigated and corroborated 7Of course, if we do not impose the additional structure of a binary allocation problem, then there exist OSP mechanisms that are not monotone price mechanisms. This paper contains several examples.
5 this claim (Kagel et al., 1987; Kagel and Levin, 1993). More generally, laboratory subjects find it difficult to reason state-by-state about hypothetical scenarios (Charness and Levin, 2009; Esponda and Vespa, 2014; Ngangoue and Weizsacker, 2015). This mental process, often called “contingent reasoning”, has received little formal treatment in economic theory.
There is also a strand of literature, including Vickrey’s seminal paper, that observes that sealed-bid auctions raise problems of commitment (Vickrey, 1961; Rothkopf et al., 1990; Cramton, 1998).
For instance, it may be difficult to prevent shill bidding without third-party verification. Rothkopf et al. (1990) argue that “robustness in the face of cheating and of fear of cheating is important in determining auction form”.
This paper formalizes and unifies both these strands of thought. It shows that mechanisms that do not require contingent reasoning are identical to mechanisms that can be run under bilateral commitment.
Eyster and Rabin (2005) and Esponda (2008) model agents who do not fully account for other agents’ private information. An extensive literature on level-k reasoning8 models agents who hold non-equilibrium beliefs about other agents’ strategies. These are conceptually distinct from mis-takes in contingent reasoning. In particular, these models predict no deviations from dominant-strategy play in strategy-proof mechanisms.9 The Prisoner’s Dilemma is a special case of game (i) in Figure 1; playing defect is not obviously dominant.
On the other hand, if Agent 1 is informed of Agent 2’s action before making his decision, then playing defect is obviously dominant. Shafir and Tversky (1992) find that laboratory subjects in a Prisoner’s Dilemma are more likely to play the weakly dominant strategy when they are informed beforehand that their opponent has cooperated (84%) or when they are informed beforehand that their opponent has defected (97%), compared to when they are not informed of their opponent’s strategy (63%).
This paper relates to the planned US auction to repurchase television broadcast rights. In this setting, complex underlying constraints have the result that Vickrey prices cannot be computed without large approximation errors. To surmount this challenge, Milgrom and Segal (2015) study a class of algorithms with desirable properties, called direct deferred-acceptance threshold auctions.
Their results imply that every such algorithm is OSP-implementable.
In combinatorial auction problems, finding the optimal solution is NP-hard, so the Vickrey-Clarke-Groves mechanism may be computationally infeasible. Consequently, there has been sub-stantial interest in ‘posted-price’ mechanisms that approximate the optimum in polynomial time (Bartal et al., 2003; Feldman et al., 2014). These have the (previously unmodeled) advantage of being obviously strategy-proof.
For some mechanisms, there exist polynomial-time algorithms that verify that the mechanism is strategy-proof (Brˆ anzei and Procaccia, 2015; Barthe et al., 2015). These are useful if agents do not trust that the mechanism is strategy-proof, but are otherwise computationally sophisticated.
OSP requires equilibrium in obviously dominant strategies. This is distinct from O-solvability, a solution concept used in the computer science literature on decentralized learning. (Friedman, 2002, 2004) Strategy Si overwhelms S′ i if the worst possible outcome from Si is strictly better than the best possible outcome from S′ i. O-solvability calls for the iterated deletion of overwhelmed strategies. One difference between the two concepts is that O-solvability is for normal form games, whereas OSP invokes a notion of an ‘earliest point of departure’, which is only defined in the extensive form.
O-solvability is too strong for our current purposes, because almost no games 8Stahl and Wilson (1994, 1995); Nagel (1995); Camerer et al. (2004); Crawford and Iriberri (2007a,b).
9Level-0 agents may deviate from dominant strategy play in a strategy-proof mechanism. However, the behavior of level-0 agents is a primitive of the theory, and a sufficiently large population of level-0 agents can explain any data.
6 studied in mechanism design are O-solvable.10 In subsequent work, Ashlagi and Gonczarowski (2015) investigate the OSP-implementation of the deferred acceptance algorithm (Gale and Shapley, 1962). Loertscher and Marx (2015) con-sider the problem of bilateral trade with many buyers and many sellers, and produce a prior-free asymptotically-optimal OSP mechanism.
3 Definition and Characterization The planner operates in an environment consisting of: 1. A set of agents, N ≡{1, . . . , n}.
2. A set of outcomes, X.
3. A set of type profiles, Θ ≡Q i∈N Θi.
4. A utility function for each agent, ui : X × Θi →R An extensive game form with consequences in X is a tuple ⟨H, ≺, A, A, P, δc, (Ii)i∈N, g⟩, where: 1. H is a set of histories, along with a binary relation ≺on H that represents precedence.
(a) ≺is a partial order, and (H, ≺) form an arborescence.
(b) h∅denotes h ∈H : ¬∃h′ : h′ ≺h (c) H has finite depth, i.e.: ∃k ∈N : ∀h ∈H : |{h′ ∈H : h′ ≺h}| ≤k (2) (d) Z ≡{h ∈H : ¬∃h′ : h ≺h′} (e) σ(h) denotes the set of immediate successors of h.
2. A is a set of actions.
3. A : H \ h∅→A labels each non-initial history with the last action taken to reach it.
(a) A is one-to-one on σ(h).
(b) A(h) denotes the actions available at h.
A(h) ≡ [ h′∈σ(h) A(h′) (3) 4. P is a player function. P : H \ Z →N ∪c 5. δc is the chance function. It specifies a probability measure over chance moves. dc denotes some realization of chance moves: For any h where P(h) = c, dc(h) ∈A(h).11 10For instance, neither ascending clock auctions nor second-price sealed-bid auctions are O-solvable.
11We could make the addition assumption that δc has full support on the available moves A(h) when it is called to play. This ensures a pleasing invariance property: It rules out games with zero-probability chance moves that do not affect play, but do affect whether a strategy is obviously dominant. However, a full support assumption is not necessary for any of the results that follow, so we do not make it here.
7 6. Ii is a partition of {h : P(h) = i} such that: (a) A(h) = A(h′) whenever h and h′ are in the same cell of the partition.
(b) For any Ii ∈Ii, we denote: P(Ii) ≡P(h) for any h ∈Ii. A(Ii) ≡A(h) for any h ∈Ii.
(c) Each action is available at only one information set: If a ∈A(Ii), a′ ∈A(I′ j), Ii ̸= I′ j then a ̸= a′.
7. g is an outcome function. It associates each terminal history with an outcome. g : Z →X Additionally we denote Ii ≺I′ i if there exist h, h′ such that: 1. h ≺h′ 2. h ∈Ii 3. h′ ∈I′ i We denote Ii ≺h if there exists h′ ∈Ii such that h′ ≺h. h ≺Ii is defined symmetrically. We use ⪯to denote the corresponding weak order.
A strategy Si for agent i in game G specifies what agent i does at every one of her information sets. Si(Ii) ∈A(Ii). A strategy profile S = (Si)i∈N is a set of strategies, one for each agent. When we want to refer to the strategies used by different types of i, we use Sθi i to denote the strategy assigned to type θi.
Let zG(h, S, δc) be the lottery over terminal histories that results in game form G when we start from h and play proceeds according to (S, δc). zG(h, S, dc) is the result of one realization of the chance moves under δc. We sometimes write this as zG(h, Si, S−i, dc).
Let uG i (h, Si, S−i, dc, θi) ≡ui(g(zG(h, Si, S−i, dc)), θi). This is the utility to agent i in game G, when we start at history h, play proceeds according to (Si, S−i, dc), and the resulting outcome is evaluated according to preferences θi.
Definition 1. ψi(h) is the experience of agent i along history h. ψi(h) is an alternating sequence of information sets and actions. It is constructed as follows: Initialize t = 1, h1 = h∅, ψi = {∅}.
1. If t > 1 and P(ht−1) = i, add A(ht) to the end of ψi.
2. If P(ht) = i, add Ii : ht ∈Ii to the end of ψi.
3. Terminate if ht = h.
4. Set ht+1 := h′ ∈H : h′ ∈σ(ht) and h′ ⪯h.
5. Set t := t + 1.
6. Go to 1.
We use Ψi to denote the set {ψi(h) : h ∈H} ∪ψ∅, where ψ∅is the empty sequence.12 An extensive game form has perfect recall if for any information set Ii, for any two histories h and h′ in Ii, ψi(h) = ψi(h′). We use ψi(Ii) to denote ψi(h) : h ∈Ii.
Definition 2. G is the set of all extensive game forms with consequences in X and perfect recall.
12Mandating the inclusion of the empty sequence has the following consequence: By looking at the set Ψi, it is not possible to infer whether P(h∅) = i.
8 A choice rule is a function f : Θ →X. If we consider stochastic choice rules, then it is a function f : Θ →∆X.13 A solution concept C is a set-valued function with domain G × Θ. It takes values in the set of strategy profiles.
Definition 3. f is C-implementable if there exists 1. G ∈G 2. ((Sθi i )θi∈Θi)i∈N such that, for all θ ∈Θ 1. (Sθi i )i∈N ∈C(G, θ).
2. f(θ) = g(zG(∅, (Sθi i )i∈N, δc)) Notice that each agent’s strategy depends just on his own type. To ease notation, we abbrebiate (Sθi i )i∈N ≡Sθ and ((Sθi i )θi∈Θi)i∈N ≡(Sθ)θ∈Θ.
Our concern is with weak implementation: We require that Sθ ∈C(G, θ), not {Sθ} = C(G, θ).
This is to preserve the analogy with canonical results for strategy-proofness, many of which assume weak implementation (Myerson, 1981; Saks and Yu, 2005).
We use “(G, (Sθ)θ∈Θ) C-implements f” to mean that (G, (Sθ)θ∈Θ) fulfils the requirements of Definition 3. We use “G C-implements f” to mean that there exists (Sθ)θ∈Θ such that (G, (Sθ)θ∈Θ) fulfils the requirements of Definition 3.
Definition 4 (Weakly Dominant). In G for agent i with preferences θi, Si is weakly dominant if: ∀S′ i : ∀S′ −i : Eδc[uG i (h∅, S′ i, S′ −i, dc, θi)] ≤Eδc[uG i (h∅, Si, S′ −i, dc, θi)] (4) Let α(Si, S′ i) be the set of earliest points of departure for Si and S′ i. That is, α(Si, S′ i) contains the information sets where Si and S′ i have made identical decisions at all prior information sets, but are making a different decision now.
Definition 5 (Earliest Points of Departure). Ii ∈α(Si, S′ i) if and only if: 1. Si(Ii) ̸= S′ i(Ii) 2. There exist h ∈Ii, S−i, dc such that h ≺zG(h∅, Si, S−i, dc).
3. There exist h ∈Ii, S−i, dc such that h ≺zG(h∅, S′ i, S−i, dc).
This definition can be extended to deal with mixed strategies14, but pure strategies suffice for our current purposes.
13For readability, we generally suppress the latter notation, but the claims that follow hold for both deterministic and stochastic choice rules. Additionally, the set X could itself be a set of lotteries. The interpretation of this is that the planner can carry out one-time public lotteries at the end of the mechanism, where the randomization is observable and verifiable.
14Three modifications are necessary: First, we change requirement 1 to be that both strategies specify different probability measures at Ii. Second, we adapt requirements 2 and 3 to hold for some realization of the mixed strategies.
Finally, we include the recursive requirement, “There does not exist I′ i ≺Ii such that I′ i ∈α(Si, S′ i).” 9 Definition 6 (Obviously Dominant). In G for agent i with preferences θi, Si is obviously dom-inant if: ∀S′ i : ∀Ii ∈α(Si, S′ i) : sup h∈Ii,S′ −i,dc uG i (h, S′ i, S′ −i, dc, θi) ≤ inf h∈Ii,S′ −i,dc uG i (h, Si, S′ −i, dc, θi) (5) In words, Si is obviously dominant if, for any deviation S′ i, conditional on reaching any earliest point of departure, the best possible outcome under S′ i is no better than the worst possible outcome under Si.15 Compare Definition 4 and Definition 6. Weak dominance is defined using h∅, the history that begins the game. Consequently, if two extensive games have the same normal form, then they have the same weakly dominant strategies. Obvious dominance is defined with histories that are in information sets that are earliest points of departure.
Thus two extensive games with the same normal form may not have the same obviously dominant strategies.16 Switching to a direct revelation mechanism may not preserve obvious dominance, so the standard revelation principle does not apply.
Weak dominance treats chance moves and other players asymmetrically.
Suppose G has a Bayes-Nash equilibrium, and we replace all the agents in N \ 1 with chance moves drawn from the Bayes-Nash equilibrium distribution. Then agent 1 has a weakly dominant strategy. By contrast, obvious dominance treats chance moves and other players symmetrically.17 Definition 7 (Strategy-Proof). S ∈SP(G, θ) if for all i, Si is weakly dominant.
Definition 8 (Obviously Strategy-Proof). S ∈OSP(G, θ) if for all i, Si is obviously dominant.
A mechanism is weakly group-strategy-proof if there does not exist a coalition that could deviate and all be strictly better offex post.
Definition 9 (Weakly Group-Strategy-Proof). S ∈WGSP(G, θ) if there does not exist a coalition ˆ N ⊆N, deviating strategies ˆ S ˆ N, non-coalition strategies S′ N\ ˆ N and chance moves dc such that: For all i ∈ˆ N: uG i (h∅, ˆ S ˆ N, S′ N\ ˆ N, dc, θi) > uG i (h∅, S ˆ N, S′ N\ ˆ N, dc, θi) (6) Obvious strategy-proofness implies weak group-strategy-proofness.
Proposition 1. If S ∈OSP(G, θ), then S ∈WGSP(G, θ).
15Obvious dominance is related to conditional dominance (Shimoji and Watson, 1998), in that both concepts evaluate a dominance relationship conditional on reaching some information set. Obvious dominance specifies that this evaluation is carried out at any earliest point of departure, using a sup-inf comparison.
16Two extensive games have the same reduced normal form if and only if they can be made identical using a small set of elementary transformations (Thompson, 1952; Elmes and Reny, 1994). Which of these transformations does not preserve obvious dominance? Elmes and Reny (1994) propose three such transformations, INT, COA, and ADD, which preserve perfect recall. Brief inspection reveals that obvious dominance is invariant under INT and COA, but varies under ADD.
17If we replace other agents with chance moves in a second-price sealed-bid auction, the remaining agent’s strategy is still not obviously dominant. This prediction is consistent with the systematic mistakes that laboratory subjects make in the Becker-DeGroot-Marschak mechanism (Becker et al., 1964; Cason and Plott, 2014).
10 Proof. We prove the contrapositive.
Suppose S / ∈WGSP(G, θ).
Then there is a coalition ˆ N that could jointly deviate to strategies ˆ S ˆ N and all be strictly better off. Fix S′ N\ ˆ N and dc such that all agents in the coalition are strictly better off. Along the resulting terminal history, there must be a first agent i in the coalition to deviate from Si to ˆ Si. That first deviation happens at some information set Ii ∈α(Si, ˆ Si).
Since agent i strictly gains from that deviation, S / ∈ OSP(G, θ).18 Corollary 1. If S ∈OSP(G, θ), then S ∈SP(G, θ).
Proposition 1 suggests a question: Is a choice rule OSP-implementable if and only if it is WGSP-implementable? Proposition 6 in Subsection 4.3 shows that this is not so.
3.1 Cognitive limitations In what sense is obvious dominance obvious? Intuitively, to see that S′ i is weakly dominated by Si, the agent must understand the entire function uG i , and check that for all opponent strategy profiles S−i, the payofffrom S′ i is no better than the payofffrom Si. By contrast, to see that S′ i is obviously dominated by Si, the agent need only know the range of the functions uG i (·, S′ i, · · · ) and uG i (·, Si, · · · ) at any earliest point of departure. Thus, obvious dominance can be recognized even if the agent has a simplified mental model of the world. We now make this point rigorously.
We define an equivalence relation between mechanisms. In words, G and G′ are i-indistinguishable if there exists a bijection from i’s information sets and actions in G onto i’s information sets and actions in G′, such that: 1. ψi is an experience in G iffψi’s image is an experience in G′.
2. Outcome x could follow experience ψi in G iffx could follow ψi’s image in G′.
Definition 10. Take any G, G′ ∈G, with information partitions Ii, I′ i and experience sets Ψi, Ψ′ i.
G and G′ are i-indistinguishable if there exists a bijection λG,G′ from Ii ∪A(Ii) to19 I′ i ∪A′(I′ i) such that: 1. ψi ∈Ψi iffλG,G′(ψi) ∈Ψ′ i 2. ∃z ∈Z : g(z) = x, ψi(z) = ψi iff∃z′ ∈Z′ : g′(z′) = x, ψ′ i(z′) = λG,G′(ψi) where we use λG,G′(ψi) to denote {λG,G′(ψk i )}T k=1, where T ∈N ∪∞.
For G and G′ that are i-indistinguishable, we define λG,G′(Si) to be the strategy that, given information set I′ i in G′, plays λG,G′(Si(λ−1 G,G′(I′ i))).
The next theorem states that obviously dominant strategies are the strategies that can be recognized as weakly dominant, by an agent who has a simplified mental model of the world.
Theorem 1. For any i, θi: Si is obviously dominant in G if and only if for every G′ that is i-indistinguishable from G, λG,G′(Si) is weakly dominant in G′.
18I thank Ilya Segal for suggesting this concise proof.
19This definition entails that λG,G′ maps Ii onto I′ i and A(Ii) onto A′(I′ i). If an information set in G was mapped onto an action in G′, then any experience involving that information set would, when passed through the bijection, result in a sequence that was not an experience, and ipso facto not an experience of G′.
11 The “if” direction permits a constructive proof. Suppose Si is not obviously dominant in G. We apply a general procedure to construct G′ that is i-indistinguishable from G, such that λG,G′(Si) is not weakly dominant. The “only if” direction proceeds as follows: Suppose there exists some G′ in the equivalence class of G, where λG,G′(Si) is not weakly dominant. There exists some earliest information set in G′ where i could gain by deviating. We then use λ−1 G,G′ to locate an information set in G, and a deviation S′ i, that do not satisfy the obvious dominance inequality. Appendix A provides the details.
One interpretation of Theorem 1 is that obviously dominant strategies are those that can be recognized as dominant given only a coarse description of the game form.
Another interpretation of Theorem 1 is that obviously dominant strategies are robust to local misunderstandings. Suppose an agent could mistake any G for any other i-indistinguishable G′.
For instance, in a second-price sealed-bid auction, he might (mistakenly) conjecture that bidding above his value increases his chance of winning with positive profits. Obviously dominant strategies are those that are robust to such mistakes.20 3.2 Supported by bilateral commitments Suppose the following extended game form ˜ G with consequences in X: As before we have a set of agents N, outcomes X, and preference profiles Q i∈N Θi. However, there is one player in addition to N: Player 0, the Planner.
The Planner has an arbitrarily rich message space M. At the start of the game, each agent i ∈N privately observes θi. Play proceeds as follows: 1. The Planner chooses one agent i ∈N and sends a query m ∈M, along with a set of acceptable replies R ⊂M 2. i observes (m, R), and chooses a reply r ∈R.
3. The Planner observes r.
4. The Planner either selects an outcome x ∈X, or chooses to send another query.
(a) If the Planner selects an outcome, the game ends.
(b) If the Planner chooses to send another query, go to Step 1.
For i ∈N, i’s strategy specifies what reply to give, as a function of his preferences, the past sequence of queries and replies between him and the Planner, and the current (m, R). That is: ˜ Si(θi, (mk, Rk, rk)t−1 k=1, mt, Rt) ∈Rt (7) We use ˜ Sθi i to denote the strategy played by type θi of agent i. Again we abbreviate ( ˜ Sθi i )i∈N ≡ ˜ Sθ N and (( ˜ Sθi i )θi∈Θi)i∈N ≡( ˜ Sθ N)θ∈Θ.
˜ S0 denotes a pure strategy for the Planner. We require that these have finite length, which ensures that payoffs are well-defined.
Definition 11. ˜ S0 is a pure strategy of finite length if there exists k ∈N such that: For all ˜ SN: ( ˜ S0, ˜ SN) results in the Planner sending k or fewer total queries.
20Obvious dominance can thus be seen as a theory of game form misconceptions, of the kind called for by Cason and Plott (2014).
12 S0 denotes the set of all pure strategies of finite length.
The standard full commitment paradigm is equivalent to allowing the Planner to commit to a unique ˜ S0 ∈S0 (or some probability measure over S0). Instead, we assume that for each agent, the Planner can commit to a subset ˆ Si 0 ⊆S0 that is measurable with respect to that agent’s observations in the game.
This is formalized as follows: Each ( ˜ S0, ˜ SN) results in some observation oi ≡(oC i , oX i ), consisting of a communication sequence between the Planner and agent i, oC i = (mk, Rk, rk)T k=1 for T ∈N, as well as some outcome oX i ∈X.21 Oi is the set of all possible observations (for agent i). We define φi : S0 × SN →Oi, where φi( ˜ S0, ˜ SN) is the unique observation resulting from ( ˜ S0, ˜ SN). Next we define, for any ˆ S0 ⊆S0: Φi(ˆ S0) ≡{oi : ∃˜ S0 ∈ˆ S0 : ∃˜ SN : oi = φi( ˜ S0, ˜ SN)} (8) For any ˆ Oi ⊆Oi: Φ−1 i ( ˆ Oi) ≡{ ˜ S0 : ∀˜ SN : φi( ˜ S0, ˜ SN) ∈ˆ Oi} (9) Definition 12. ˆ S0 is i-measurable if there exists ˆ Oi such that: ˆ S0 = Φ−1 i ( ˆ Oi) (10) Intuitively, the i-measurable subsets of S0 are those such that, if the Planner deviates, then there exists an agent strategy profile such that agent i detects the deviation. Formally, the i-measurable subsets of S0 are the σ-algebra generated by Φi (where we impose the discrete σ-algebra on Oi).
Definition 13. A mixed strategy of finite length over ˆ S0 specifies a probability measure over a subset ¯ S0 ⊆ˆ S0 such that: There exists k ∈N such that: For all ˜ S0 ∈¯ S0 and all ˜ SN: ( ˜ S0, ˜ SN) results in the Planner sending k or fewer total queries.
We use ∆ˆ S0 to denote the mixed strategies of finite length over ˆ S0. ˜ S∆ 0 denotes an element of such a set.
Definition 14. A choice rule f is supported by bilateral commitments (ˆ Si 0)i∈N if 1. For all i ∈N: ˆ Si 0 is i-measurable.
2. There exist ˜ S∆ 0 , and ( ˜ Sθ N)θ∈Θ such that: (a) For all θ: ( ˜ S∆ 0 , ˜ Sθ N) results in f(θ).
(b) ˜ S∆ 0 ∈∆T i∈N ˆ Si 0 (c) For all i ∈N, θi, ˜ S′ N\i, ˜ S∆′ 0 ∈∆ˆ Si 0: ˜ Sθi i is a best response to ( ˜ S∆′ 0 , ˜ S′ N\i) given preferences θi.
Requirement 2.a is that the Planner’s mixed strategy and the agent’s pure strategies result in the (distribution over) outcomes required by the choice rule. Requirement 2.b is that the Planner’s strategy is a (possibly degenerate) mixture over pure strategies compatible with every bilateral commitment. 2.c is that each agent i’s assigned strategy is weakly dominant, when we consider the Planner as a player restricted to playing mixtures over strategies in ˆ Si 0.
21The communication sequence might be empty, which we represent using T = 0.
13 “Supported by bilateral commitments” is just one of many partial commitment regimes. This one requires that the commitment offered to each agent is measurable with respect to events that he can observe. In reality, contracts are seldom enforceable unless each party can observe breaches.
Thus, “supported by bilateral commitments” is a natural case to study.
Theorem 2. f is OSP-implementable if and only if there exist bilateral commitments (ˆ Si 0)i∈N that support f.
The intuition behind the proof is as follows: A bilateral commitment ˆ Si 0 is essentially equivalent to the Planner committing to ‘run’ only games in some i-indistinguishable equivalence class of G.
Consequently, we can find a set of bilateral commitments that support fy if and only if we can find some (G, (Sθ)θ∈Θ) such that, for every i, for every θi, for every G′ that is i-indistinguishable from G, λG,G′(Sθi i ) is weakly dominant in G′. By Theorem 1, this holds if and only if fy is OSP-implementable. Appendix A provides the details.
3.3 The pruning principle The standard revelation principle does not hold for OSP mechanisms; converting an OSP mech-anism into the corresponding direct revelation mechanism may not preserve obvious dominance.
However, there is a weaker principle that substantially simplifies the analysis.22 Here we define the pruning of a mechanism with respect to a set of strategies (one for each type of each agent). This is the new mechanism constructed by deleting all sub-trees that are not reached given any type profile.
Definition 15 (Pruning). Take any G = ⟨H, ≺, A, A, P, δc, (Ii)i∈N, g⟩, and (Sθ)θ∈Θ. P(G, (Sθ)θ∈Θ) ≡ ⟨˜ H, ˜ ≺, ˜ A, ˜ A, ˜ P, ˜ δc, (˜ Ii)i∈N, ˜ g⟩is the pruning of G with respect to (Sθ)θ∈Θ, constructed as follows: 1. ˜ H = {h ∈H : ∃θ : ∃dc : h ⪯zG(∅, Sθ, dc)} 2. For all i, if Ii ∈Ii then (Ii ∩˜ H) ∈˜ Ii.23 3. (˜ ≺, ˜ A, ˜ A, ˜ P, ˜ δc, ˜ g) are (≺, A, A, P, δc, g) restricted to ˜ H.
It turns out that, if some mechanism OSP-implements a choice rule, then the pruning of that mechanism with respect to the equilibrium strategies OSP-implements that same choice rule. Thus, while we cannot restrict our attention to direct revelation mechanisms, we can restrict our attention to ‘minimal’ mechanisms, where no histories are offthe path of play.
Proposition 2 (The Pruning Principle). Let ˜ G ≡P(G, (Sθ)θ∈Θ), and ( ˜ Sθ)θ∈Θ be (Sθ)θ∈Θ restricted to ˜ G. If (G, (Sθ)θ∈Θ) OSP-implements f, then ( ˜ G, ( ˜ Sθ)θ∈Θ) OSP-implements f.
4 Applications 4.1 Binary Allocation Problems We now consider a canonical environment, (N, X, Θ, (ui)i∈N). Let Y ⊆2N be the set of feasible allocations, with representative element y ∈Y . An outcome consists of an allocation y ∈Y and a transfer for each agent, X = Y × Rn. t ≡(ti)i∈N denotes a profile of transfers.
22For instance, the pruning principle is required to concisely state Theorem 3.
23Note that the empty history h∅is distinct from the empty set. That is to say, (Ii ∩˜ H) = ∅does not entail that {h∅} ∈˜ Ii.
14 Preferences are quasilinear. Θ = Q i∈N Θi, where Θi = [θi, θi], for 0 ≤θi < θi < ∞. For θi ∈Θi ui(θi, y, t) = 1i∈yθi + ti (11) For instance, in a private value auction with unit demand, i ∈y if and only if agent i receives at least one unit of the good under allocation y. In a procurement auction, i ∈y if and only if i does not incur costs of provision under allocation y. θi is agent i’s cost of provision (equivalently, benefit of non-provision). In a binary public goods game, Y = {∅, N}.
An allocation rule is a function fy : Θ →Y . A choice rule is thus a combination of an allocation rule and a payment rule, f = (fy, ft), where ft : Θ →Rn. Similarly, for each game form G, we disaggregate the outcome function, g = (gy, gt).
In this part, we concern ourselves only with deterministic allocation rules and payment rules, and thus suppress notation involving δc and dc.
Definition 16. An allocation rule fy is C-implementable if there exists ft such that (fy, ft) is C-implementable. G C-implements fy if there exists ft such that G C-implements (fy, ft) Definition 17. fy is monotone if for all i, for all θ−i, 1i∈fy(θi,θ−i) is weakly increasing in θi.
In this environment, fy is SP-implementable if and only if fy is monotone. This result is implicit in Spence (1974) and Mirrlees (1971), and is proved explicitly in Myerson (1981).24 Moreover, if an allocation rule fy is SP-implementable, then the accompanying payment rule ft is essentially unique.
ft,i(θi, θ−i) = −1i∈fy(θ) inf{θ′ i : i ∈fy(θ′ i, θ−i)} + ri(θ−i) (12) where ri is some arbitrary deterministic function of the other agents’ preferences. This follows easily by arguments similar to those in Green and Laffont (1977) and Holmstr¨ om (1979).
We are interested in how these results change when we require OSP-implementation. In partic-ular: 1. What condition on fy characterizes the set of OSP-implementable allocation rules?
2. For OSP-implementation, is there an analogous ‘essential uniqueness’ result on the extensive game form G?
We now define a monotone price mechanism. Informally, a monotone price mechanism is such that, for every i, starting from any information set where i has more than one action: 1. Either: (a) There is a ‘going transfer’ associated with being in the allocation, which falls monoton-ically.
(b) Whenever the going transfer falls, i chooses whether to keep bidding or to quit.
(c) If i quits, then i is not in the allocation and receives a fixed transfer (d) If the game ends: i. If i is in the allocation, then i receives the going transfer.
24These monotonicity results for are for weak SP-implementation rather than full SP-implementation implemen-tation. Weak SP-implementation requires Sθ ∈SP(G, θ). Full SP-implementation requires {Sθ} = SP(G, θ). There are monotone allocation rules for which the latter requirement cannot be satisfied. For example, suppose two agents with unit demand. Agent 1 receives one unit iffv1 > .5. Agent 2 receives one unit iffv2 > v1.
15 ii. If i is not in the allocation, then i receives the fixed transfer.
(e) If i has more than one non-quitting action: i. The going transfer will not fall in future.
ii. An action is available that guarantees i is in the allocation.
2. Or: (a) There is a ‘going transfer’ associated with not being in the allocation, which falls mono-tonically.
(b) Whenever the going transfer falls, i chooses whether to keep bidding or to quit.
(c) If i quits, then i is in the allocation and receives a fixed transfer (d) If the game ends: i. If i is not in the allocation, then i receives the going transfer.
ii. If i is in the allocation, then i receives the fixed transfer.
(e) If i has more than one non-quitting action: i. The going transfer will not fall in future.
ii. An action is available that guarantees i is not in the allocation.
The “either” clause contains ascending clock auctions as a special case. The “or” clause contains descending price procurement auctions; agents that do not win the contract receive a fixed zero transfer. There is a positive payment associated with winning the contract (i.e. not being in the allocation), which starts high and counts downwards.
The following definition restates the above formally.
Definition 18 (Monotone Price Mechanism). A game G is a monotone price mechanism if, for every i ∈N, at every earliest information set I∗ i such that A(I∗ i ) > 1: 1. Either: There exists a real number t0 i , a function ˜ t1 i : {Ii : I∗ i ∈ψi(Ii)} →R, and a set of actions A0 such that: (a) For all a ∈A0, for all z such that a ∈ψi(z): i / ∈gy(z) and gt,i(z) = t0 i .
(b) A0 ∩A(I∗ i ) ̸= ∅.
(c) For all I′ i, I′′ i ∈{Ii : I∗ i ⪯Ii}: i. If I′ i ≺I′′ i , then ˜ t1 i (I′ i) ≥˜ t1 i (I′′ i ).
ii. If I′ i is the penultimate information set in ψi(I′′ i ) and ˜ t1 i (I′ i) > ˜ t1 i (I′′ i ), then A0 ∩ A(I′′ i ) ̸= ∅.
iii. If I′ i ≺I′′ i and ˜ t1 i (I′ i) > ˜ t1 i (I′′ i ), then |A(I′ i) \ A0| = 1.
iv. If |A(I′ i) \ A0| > 1, then there exists a ∈A(I′ i) such that: For all z such that a ∈ψi(z): i ∈gy(z).
(d) For all z where I∗ i ≺z: i. Either: i / ∈gy(z) and gt,i(z) = t0 i .
ii. Or: i ∈gy(z) and gt,i(z) = inf Ii:I∗ i ⪯Ii≺z ˜ t1 i (Ii) (13) 16 2. Or: There exists a real number t1 i , a function ˜ t0 i : {Ii : I∗ i ∈ψi(Ii)} →R, and a set of actions A1 such that: (a) For all a ∈A1, for all z such that a ∈ψi(z): i ∈gy(z) and gt,i(z) = t1 i .
(b) A1 ∩A(I∗ i ) ̸= ∅.
(c) For all I′ i, I′′ i ∈{Ii : I∗ i ⪯Ii}: i. If I′ i ≺I′′ i , then ˜ t0 i (I′ i) ≥˜ t0 i (I′′ i ).
ii. If I′ i is the penultimate information set in ψi(I′′ i ) and ˜ t1 i (I′ i) > ˜ t1 i (I′′ i ), then A1 ∩ A(I′′ i ) ̸= ∅.
iii. If I′ i ≺I′′ i and ˜ t0 i (I′ i) > ˜ t0 i (I′′ i ), then |A(I′ i) \ A1| = 1.
iv. If |A(I′ i) \ A1| > 1, then there exists a ∈A(I′ i) such that: For all z such that a ∈ψi(z): i / ∈gy(z).
(d) For all z where I∗ i ≺z: i. Either: i ∈gy(z) and gt,i(z) = t1 i .
ii. Or: i / ∈gy(z) and gt,i(z) = inf Ii:I∗ i ⪯Ii≺z ˜ t0 i (Ii) (14) Notice what this definition does not require. The going transfer need not be equal across agents.
Whether and how much one agent’s going transfer changes could depend on other agents’ actions.
Some agents could face a procedure consistent with the ‘either’ clause, and other agents could face a procedure consistent with the ‘or’ clause. Indeed, which procedure an agent faces could depend on other agents’ actions.
Theorem 3. If (G, (Sθ)θ∈Θ) OSP-implements fy, then ˜ G ≡P(G, (Sθ)θ∈Θ) is a monotone price mechanism.
Ascending auctions are characterized by a number of fine procedural details: Each agent faces a ‘going price’. That price only ever rises. At each point, the agent chooses whether to keep bidding or to quit. If an agent quits, he quits irrevocably. If the agent wins, he pays the going price. By Theorem 3, obvious dominance entails many of these details, and can be seen as an answer to the question, “Why are ascending auctions so common?” Theorem 4 shows that Theorem 3 is a full characterization.
Theorem 4. If G is a monotone price mechanism, then there exists fy such that G OSP-implements fy.
To characterize the set of OSP-implementable allocation rules, we define traversable allocation rules. We build towards this definition step by step. All vector inequalities are in the product order.25 Consider the set: {θ : ∀θ′ N\i ≥θN\i : i / ∈fy(θi, θ′ N\i)} (15) These are the type profiles such that, if we hold i’s type constant and (weakly) increase every other agent’s type, then fy requires that i is not in the allocation. To build intuition, suppose we are running an ascending clock auction with agent-specific clocks. No agent has quit yet, and the 25That is, v ≥v′ ifffor every index i, vi ≥v′ i. Similarly, v > v′ ifffor every index i, vi > v′ i.
17 current price vector is in the set defined by Equation 15. Then we can safely advance i’s clock while respecting the allocation rule.
Suppose instead that only a subset of the agents B ⊆N are actively bidding, and we already know the types of the other agents θN\B. Then Equation 16 is a natural generalization of Equation 15.
{θB : ∀θ′ B\i ≥θB\i : i / ∈fy(θi, θ′ B\i, θN\B)} (16) Definition 19. fy is traversable if: For all B ⊆N, for all θN\B, for ˜ ΘB(θN\B) ≡∩i∈Bclosure({θB : ∀θ′ B\i ≥θB\i : i / ∈fy(θi, θ′ B\i, θN\B)}) (17) 1. ˜ ΘB(θN\B) is connected.
2. There exists i ∈B such that, if θB > sup{˜ ΘB(θN\B)}, then i ∈fy(θB, θN\B).
The sets defined by Equation 17 are join-semilattices.26 There exist traversable allocation rules that are not monotone, and vice versa.27 The next theorem characterizes the set of OSP-implementable allocation rules. It invokes two additional assumptions.
First, we assume that fy admits a finite partition, which means that we can partition the type space into a finite set of |N|-dimensional intervals, with the allocation rule constant within each interval. This assumption is largely technical. It is required because OSP is defined for extensive game forms such that play proceeds in discrete steps.
OSP is not defined for continuous-time auctions, although we can approximate some of them arbitrarily finely.28 Definition 20. fy admits a finite partition if there exists K ∈N such that, for each i, there exists {θk i }K k=1 such that: 1. θi = θ1 i < θ2 i < . . . < θK i = θi.
2. For all θi, θ′ i, for all θ−i, if there does not exist k such that θi ≤θk i < θ′ i, then fy(θi, θ−i) = fy(θ′ i, θ−i) Second, we assume that the lowest type of each agent is never in the allocation, and has a zero transfer. This is a substantive restriction, and rules out, for instance, the subsidized-trade mechanisms examined by Myerson and Satterthwaite (1983).
Theorem 5. Assume that: 1. fy admits a finite partition.
2. For all i, for all θ−i, i / ∈fy(θi, θ−i).
There exists G and ft such that: 26For a proof, see Lemma 6 in Appendix A.
27Let |N| = 2 and Θ1 = Θ2 = [0, 1]. The following allocation rule is traversable but not monotone: i ∈fy(θi, θ−i) if and only if θi > θ−i and θi ≤1 2. The following rule is monotone but not traversable: i ∈fy(θi, θ−i) if and only if θi > 1 3 and θ−i ≤2 3.
28Simon and Stinchcombe (1989) show that discrete time with a very fine grid can be a good proxy for continuous time. However, in their theory, players have perfect information about past activity in the system. Adapting this to our theory, where G includes all discrete-time game forms with imperfect information, is far from straightforward.
18 1. G OSP-implements (fy, ft) 2. For all i, for all θ−i, ft,i(θi, θ−i) = 0 if and only if fy is monotone and traversable.
4.2 Online Advertising Auctions We now study an online advertising environment, which generalizes Edelman et al. (2007).
There are n bidders, and n −1 advertising positions.29 Each position has an associated click-through rate αk, where α1 ≥α2 ≥. . . ≥αn−1 > 0. For convenience, we define position n with αn = 0.
Each bidder’s type is a vector, θi ≡(θk i )n k=1. A bidder with type θi who receives position k and transfer t has utility: ui(k, t, θi) = αkθk i + t (18) The marginal utility of moving to position k from position k′, for type θi, is m(k, k′, θi) ≡αkθk i −αk′θk′ i (19) We make the following assumptions on the type space Θ: A1. Finite: |Θ| < ∞ (20) A2. Higher slots are better: ∀k ≤n −1 : ∀θ ∈Θ : ∀i ∈N : m(k, k + 1, θi) ≥0 (21) A3. Single-crossing:30 ∀k ≤n −2 : ∀θ, θ′ ∈Θ : ∀i, j ∈N : If m(k, k + 1, θi) > m(k, k + 1, θ′ j), then m(k + 1, k + 2, θi) > m(k + 1, k + 2, θ′ j).
(22) A1 is a technical assumption to accommodate extensive game forms that move in discrete steps.
A2 and A3 are substantive assumptions. Edelman et al. (2007) assume that for all k, k′, θk i = θk′ i , which entails A2. If α1 > α2 > . . . > αn−1 > 0, then their assumption also entails A3.
In this environment, the Vickrey-Clarke-Groves (VCG) mechanism selects the efficient alloca-tion.
Suppose we number each buyer according to the slot he wins.
Then bidder i has VCG payment: −ti = n−1 X k=i m(k, k + 1, θk+1) (23) Edelman et al. (2007) produce a generalized English auction that ex post implements the efficient allocation rule in online advertising auctions. The generalized English auction has a unique perfect Bayesian equilibrium in continuous strategies. It is not SP, and therefore is not OSP.
Here we produce an alternative ascending auction that OSP-implements the efficient allocation rule.
29It is trivial to extend what follows to fewer than n −1 advertising positions, but doing so would add notation.
30This assumption is not identical to the single-crossing assumption in Yenmez (2014). For instance, Yenmez’s condition permits the second inequality in Equation 22 to be weak.
19 Proposition 3. Assume A1, A2, A3. There exists G that OSP-implements the efficient allocation rule and the VCG payments.
Proof. We construct G. Set pn−1 := 0, An−1 = N.
For l = 1, . . . , n −1: 1. Start the price at pn−l.
2. Raise the price in small increments. If the current price is p′ n−l, the next price is: p′′ n−l := inf θ∈Θ,i∈N{m(n −l, n −l + 1, θi) : m(n −l, n −l + 1, θi) > p′ n−l} (24) 3. At each price, query each agent in An−l (in an arbitrary order), giving her the option to quit.
4. At any price p′ n−l, if agent i quits, allocate her slot n−l+1, and charge every agent in An−l \i the price p′ n−l.
5. Set pn−l−1 := inf θ∈Θ,i∈N{m(n −l −1, n −l, θi) : m(n −l, n −l + 1, θi) ≥p′ n−l} (25) An−l−1 := An−l \ i (26) It is an obviously dominant strategy for agent i to quit iffthe price in round l is weakly greater than m(n −l, n −l + 1, θi).
Consider any round l. Payments from previous rounds are sunk costs. Quitting yields slot n −l + 1 at no additional cost, and removes the agent from future rounds.
Consider deviations where the earliest point of departure involves quitting. The current price p′ n−l is weakly less than m(n −l, n −l + 1, θi). If the truth-telling strategy has the result that i quits in round l, this outcome is at least as good for i as quitting now. If the truth-telling strategy has the result that i does not quit in round l, then i is charged some amount less than his marginal value for moving up a slot, and the next starting price is pn−l−1 ≤m(n −l −1, n −l, θi), so the argument repeats.
Consider deviations where the earliest point of departure involves staying in. The current price p′ n−l is weakly greater than m(n −l, n −l + 1, θi), so this either has the same result as quitting now, or raises i’s position at marginal cost weakly above i’s marginal utility. This is trivially true for the current round.
Consider the next round, l + 1.
If the starting price pn−l−1 is strictly less m(n −l −1, n −l, θi), then there exists some θ′ and j such that m(n −l −1, n −l, θi) > m(n −l −1, n −l, θ′ j). And m(n −l, n −l + 1, θi) ≤p′ n−l ≤m(n −l, n −l + 1, θ′ j), which contradicts A3. Repeating the argument suffices to prove the claim for all rounds l′ ≥l.
By inspection, this mechanism and the specified strategy profile result in the efficient allocation and the VCG payments.
Internet transactions conducted by a central auctioneer raise commitment problems, and bidders may be legitimately concerned about shill bidding. If we consider such auctions as repeated games, 20 reputation can ameliorate commitment problems, but the set of equilibria can be very large and prevent tractable analysis.
Proposition 3 shows that, even if we do not consider such auctions as repeated games, there sometimes exist robust mechanisms that rely only on bilateral commitments. In the case of adver-tising auctions, the speed of transactions may require bidders to implement their strategies using automata.
4.3 Top Trading Cycles We now produce an impossibility result in a classic matching environment (Shapley and Scarf, 1974).
There are n agents in the market, each endowed with an indivisible good. An agent’s type is a vector θi ∈Rn. Θ is the set of all n by n matrices of real numbers. An outcome assigns one object to each agent. If agent i is assigned object k, he has utility θk i . There are no money transfers.
Given preferences θ and agents R ⊆N, a top trading cycle is a set ∅⊂R′ ⊆R whose members can be indexed in a cyclic order: R′ = {i1, i2, . . . , ir = i0} (27) such that each agent ik likes ik+1’s good at least as much as any other good in R. Following Roth (1982), we assume that the algorithm in question has an arbitrary, fixed way of resolving ties.
Definition 21. f is a top trading cycle rule if, for all θ, f(θ) is equal to the output of the following algorithm: 1. Set R1 := N 2. For l = 1, 2, , . . .: (a) Choose some top trading cycle R′ ⊆Rl.
(b) Carry out the indicated trades.
(c) Set Rl+1 := Rl \ R′.
(d) Terminate if Rl+1 = ∅.
The above algorithm is of economic interest, because it finds a core allocation in an economy with indivisible goods (Shapley and Scarf, 1974).
Proposition 4. If f is a top trading cycle rule, then there exists G that SP-implements f. (Roth, 1982) Proposition 5. If f is a top trading cycle rule, then there exists G that WGSP-implements f.
(Bird, 1984) Proposition 6. If f is a top trading cycle rule and n ≥3, then there does not exist G that OSP-implements f.
Proof. SP-implementability is a hereditary property of functions. That is, if f is SP-implementable given domain Θ, then the subfunction f′ = f with domain Θ′ ⊆Θ is SP-implementable.
By inspection, the same is true for OSP-implementability. Thus, to prove Proposition 6, it suffices to produce a subfunction that is not OSP-implementable.
21 Consider the following subset Θ′ ⊂Θ. Take agents a, b, c, with endowed goods A, B, C. a has only two possible types, θa and θ′ a, such that Either B ≻a C ≻a A ≻a . . .
or C ≻a B ≻a A ≻a . . .
(28) We make the symmetric assumption for b and c.
We now argue by contradiction.
Take any G pruned with respect to the truthful strategy profiles, such that (by Proposition 2) G OSP-implements f′ = f for domain Θ′. Consider some history h at which P(h) = a with a non-singleton action set. This cannot come before all such histories for b and c.
Suppose not, and suppose B ≻a C. If a chooses the action corresponding to B ≻a C, and faces opponent strategies corresponding to C ≻b A and B ≻c A, then a receives good A. If a chooses the action corresponding to C ≻a B, and faces opponent strategies corresponding A ≻c B, then a receives good C. Thus, it is not an obviously dominant strategy to choose the action corresponding to B ≻a C. So a cannot be the first to have a non-singleton action set.
By symmetry, this argument applies to b and c as well. So all of the action sets for a, b, and c are singletons, and G does not OSP-implement f′, a contradiction.
Proposition 6 implies that the OSP-implementable choice rules are not identical to the WGSP-implementable choice rules.
5 Laboratory Experiment Are obviously strategy-proof mechanisms easier for real people to understand?
The following laboratory experiment provides a straightforward test: We compare pairs of mechanisms that implement the same allocation rule. One mechanism in each pair is SP, but not OSP. The other mechanism is OSP. Standard game theory predicts that both mechanisms will produce the same outcome. We are interested in whether subjects play the dominant strategy at higher rates under OSP mechanisms.
5.1 Experiment Design The experiment is an across-subjects design, comparing three pairs of games. There are four players in each game.
For the first pair, we compare the second-price auction (2P) and the ascending clock auction (AC). In both these games, subjects bid for a money prize. Subjects have induced affiliated private values; if a subject wins the prize, he earns an amount equal to the value of the prize, minus his payments from the auction. For each subject, his value for the prize is equal to a group draw plus a private adjustment. The group draw is uniformly distributed between $10 and $110. The private adjustment is uniformly distributed between $0 and $20. All money amounts in these games are in 25-cent increments. Each subject knows his own value, but not the group draw or the private adjustment.31 31We use affiliated private values for two reasons. First, in strategy-proof auctions with independent private values, incentives for truthful bidding are weak for bidders with values near the extremes. Affiliation strengthens incentives for these bidders. Second, Kagel et al. (1987) use affiliated private values, and the first part of the experiment is designed to replicate their results.
22 2P is SP, but not OSP. In 2P, subjects submit their bids simultaneously. The highest bidder wins the prize, and makes a payment equal to the second-highest bid. Bids are constrained to be between $0 and $150.32 AC is OSP. In AC, the price starts at a low value (the highest $25 increment that is below the group draw), and counts upwards, up to a maximum of $150. Each bidder can quit at any point.33 When only one bidder is left, that bidder wins the object at the current price.
Previous studies comparing second-price auctions to ascending clock auctions have small sam-ple sizes, given that when the same subjects play a sequence of auctions, these are plainly not independent observations. Kagel et al. (1987) compare 2 groups playing second-price auctions to 2 groups playing ascending clock auctions. Harstad (2000) compares 5 groups playing second-price auctions to 3 groups playing ascending clock auctions. (The comparison is not the main goal of either experiment.) Other studies find similar results for second-price auctions (Kagel and Levin, 1993) and for ascending clock auctions (McCabe et al., 1990), but these do not directly compare the two formats with the same value distribution and the same subject pool. When we compare 2P and AC, we can see this as a high-powered replication of Kagel et al. (1987), since we now observe 18 groups playing 2P and 18 groups playing AC.34 For the second pair, we compare the second-price plus-X auction (2P+X) and the ascending clock plus-X auction (AC+X). Subjects’ values are drawn as before. However, there is an addi-tional random variable X, which is uniformly distributed between $0 and $3. Subjects are not told the value of X until after the auction.
2P+X is SP, but not OSP. In 2P+X, subjects submit their bids simultaneously. The highest bidder wins the prize if and only if his bid exceeds the second-highest bid plus X. If the highest bidder wins the prize, then he makes a payment equal to the second-highest bid plus X. Otherwise, no agent wins the prize, and no payments are made. In this game, it is a dominant strategy to submit a bid equal to your value.
AC+X is OSP. In AC+X, the price starts at a low value (the highest $25 increment that is below the group draw), and counts upwards. Each bidder can quit at any point.35 When only one bidder is left, the price continues to rise for another X dollars, and then freezes. If the highest bidder keeps bidding until the price freezes, then she wins the prize at the final price. Otherwise, no agent wins the prize and no payments are made. In this game, it is an obviously dominant strategy to keep bidding if the price is strictly below your value, and quit otherwise.
Some subjects might find 2P or AC familiar, since such mechanisms occur in some natural economic environments. Differences in subject behavior might be caused by different degrees of familiarity with the mechanism.
2P+X and AC+X are novel mechanisms that subjects are unlikely to find familiar. 2P+X and AC+X can be seen as perturbations of 2P and AC; the underlying allocation rule is made more complex while preserving the SP-OSP distinction. Thus, comparing 2P+X and AC+X indicates whether the distinction between SP and OSP mechanisms holds for novel and more complicated auction formats.
In the third pair of games, subjects may receive one of four common-value money prizes. The four prize values are drawn, uniformly at random and without replacement, from the set: {$0.00, $0.25, $0.50, $0.75, $1.00, $1.25} (29) 32In both 2P and AC, if there is a tie for the highest bid, then no bidder wins the object.
33During the auction, each bidder observes the number of active bidders.
34I am not aware of any previous laboratory experiment that directly compares second-price and ascending clock auctions, holding constant the value distribution and subject pool, with more than five groups playing each format.
35As in AC, each bidder observes the number of active bidders. However, if the number of active bidders is 1 or 2, then the computer display informs bidders that the number of active bidders is “1 or 2”.
23 Table 1: Mechanisms in each treatment 10 rounds 10 rounds 10 rounds Treatment 1 AC AC+X OSP-RSD Treatment 2 2P 2P+X SP-RSD Treatment 3 AC AC+X SP-RSD Treatment 4 2P 2P+X OSP-RSD Subjects observe the values of all four prizes at the start of each game.
In a strategy-proof random serial dictatorship (SP-RSD), subjects are informed of their priority score, which is drawn uniformly at random from the integers 1 to 10. They then simultaneously submit ranked lists of the four prizes. Players are processed sequentially, from the highest priority score to the lowest. Ties in priority score are broken randomly. Each player is assigned the highest-ranked prize on his list, among the prizes that have not yet been assigned. It is a dominant strategy to rank the prizes in order of their money value. SP-RSD is SP, but not OSP.
In an obviously strategy-proof random serial dictatorship (OSP-RSD), subjects are informed of their priority score. Players take turns, from the highest priority score to the lowest. When a player takes his turn, he is shown the prizes that have not yet been taken, and picks one of them.
It is an obviously dominant strategy to pick the available prize with the highest money value.
SP-RSD and OSP-RSD differ from the auctions in several ways. The auctions are private-value games of incomplete information, whereas SP-RSD and OSP-RSD are common-value games of complete information. In the auctions, subjects face two sources of strategic uncertainty: They are uncertain about their opponents’ valuations, and they are uncertain about their opponents’ strategies (a function of valuations). By contrast, in SP-RSD and OSP-RSD, subjects face no uncertainty about their opponents’ valuations.
Unlike the auctions, SP-RSD and OSP-RSD are constant-sum games, such that one player’s action cannot affect total player surplus. Any effect that persists in both the auctions and the serial dictatorships is difficult to explain using social preferences, since such theories typically make different predictions for constant-sum and non-constant-sum games. Thus, in comparing SP-RSD and OSP-RSD, we test whether the SP-OSP distinction has empirical support in mechanisms that are very different from auctions.
At the start of the experiment, subjects are randomly assigned into groups of four.
These groups persist throughout the experiment. Consequently, each group’s play can be regarded as a single independent observation in the statistical analysis.36 Each group either plays 10 rounds of AC, followed by 10 rounds of AC+X, or plays 10 rounds of 2P, followed by 10 rounds of 2P+X.37 At the end of each round, subjects are shown the auction result, their own profit from this round, the winning bidder’s profit from this round, and the bids (in order from highest to lowest). Notice that subjects have 10 rounds of experience with a standard auction, before being presented with its unusual +X variant. Thus, the data from +X auctions 36We could, alternatively, randomly rematch subjects after each round. However, in order for subjects to avoid encountering past opponents, this would require more than 50 subjects per session, which is more than can be accommodated in a typical experimental economics laboratory. Moreover, we would have to cluster standard errors at the session level, requiring many sessions with more than 50 subjects per session.
This violates our budget constraint.
37If a stage game with dominant strategies is repeated finitely many times, then the resulting repeated game typically does not have a dominant strategy.
The same holds for obviously dominant strategies.
Consequently, in interpreting these results as informing us about dominant strategy play, we invoke an implicit narrow framing assumption. The same assumption is made for other experiments in this literature, such as Kagel et al. (1987) and Kagel and Levin (1993).
24 record moderately experienced bidders grappling with a new auction format.
Next, groups are re-randomized into either 10 rounds of OSP-RSD or 10 rounds of SP-RSD.
At the end of each round, subjects see which prize they have obtained, and whether their priority score was the highest, or second-highest, and so on.
Table 1 summarizes the design. Subjects had printed copies of the instructions, and the exper-imenter read aloud the part pertaining to each 10-round segment just before that segment began.
The instructions (correctly) informed subjects that their play in earlier segments would not affect the games in later segments. The instructions did not mention dominant strategies or provide recommendations for how to play, so as to prevent confounds from the experimenter demand effect.
Instructions for both SP and OSP mechanisms are of similar length and similar reading levels38, and can be found in Appendix D.
In every SP mechanism, each subject had 90 seconds to make his choice. Each subject could revise his choice as many times as he desired during the 90 seconds, and only his final choice would count. For OSP mechanisms, mean time to completion was 113.0 seconds in AC, 121.4 seconds in AC+X, and 40.5 seconds in RSD-OSP. However, the rules of the OSP mechanisms imply that not every subject was actively choosing throughout that time.
5.2 Administrative details Subjects were paid $20 for participating, in addition to their profits or losses from every round of the experiment. On average, subjects made $37.54, including the participation payment. Subjects who made negative net profits received just the $20 participation payment.
I conducted the experiment at the Ohio State University Experimental Economics Laboratory in August 2015, using z-Tree (Fischbacher, 2007). I recruited subjects from the student population using an online system. I administered 16 sessions, where each session involved 1 to 3 groups. Each session lasted about 90 minutes. In total, the data include 144 subjects in 36 groups of 4 (with 9 groups in each treatment).39 54% of subjects are male, and 21% self-report as being economics majors.
5.3 Statistical Analysis The data include 4 different auction formats, with 180 auctions per format, for a total of 720 auctions.40 One natural summary statistic for each auction is the difference between the second-highest bid and the second-highest value. This is, equivalently, the difference between that auction’s closing price, and the closing price that would have occurred if all bidders played the dominant strategy.
Figure 2 displays histograms of the second-highest bid minus the second-highest value, for AC and 2P. Figure 3 does the same for AC+X and 2P+X. If all agents are playing the dominant strategy in an auction, then the histogram for that auction will be a point mass at zero.
There is a substantial difference between the empirical distributions for OSP and SP mecha-nisms. If we choose a random auction from the data, how likely is it to have a closing price within $2.00 of the dominant strategy price? An auction is 31 percentage points more likely to have a closing price within $2.00 of the dominant strategy price under AC (OSP) compared to 2P (SP).
38Both sets of instructions are approximately at a fifth-grade reading level according to the Flesch-Kincaid read-ability test, which is a standard measure for how difficult a piece of text is to read (Kincaid et al., 1975).
39In two cases, network errors caused crashes which prevented a group from continuing in the experiment.
I recruited new subjects to replace these groups.
40In 2 out of 720 auctions, computer errors prevented bidders from correctly entering their bids. We omit these 2 observations, but including them has little effect on any of the results that follow.
25 Figure 2: Histogram: 2nd-highest bid minus 2nd-highest value for AC and 2P An auction is 28 percentage points more likely to have a closing price within $2.00 of the dominant strategy price under AC+X (OSP) compared to 2P+X (SP). Closing prices under 2P+X are systematically biased upwards (p = .0031)41.
Table 2 displays the mean absolute difference between the second-highest bid and the second-highest value, for the first 5 rounds and the last 5 rounds of each auction. This measures the magnitude of errors under each mechanism. (Alternative measures of errors are in Appendix B.) Errors are systematically larger under SP than under OSP, and this difference is significant in both the standard auctions and the novel +X auctions, and in both early and late rounds. To build intuition for effect sizes, consider that the expected profit of the winning bidder in 2P and AC is about $4.00 (given dominant strategy play). Thus, the average errors under 2P are larger than the theoretical prediction for total bidder surplus.
There is some evidence of learning in 2P; errors are smaller in the last five rounds compared to the first five rounds (p = .045, paired t-test). For the other three auction formats, there is no significant evidence of learning.42 To compare subject behavior under SP-RSD and OSP-RSD, we compute the proportion of games that do not end in the dominant strategy outcome. Under SP-RSD, 36.1% of games do not end in the dominant strategy outcome. Under OSP-RSD, 7.2% of games do not end in the dominant strategy outcome. Table 3 displays the empirical frequency of non-dominant strategy outcomes, by format and by 5-round blocks.
Deviations from the dominant strategy outcome happen more frequently under SP-RSD than under OSP-RSD, and these differences are highly significant in both early and late rounds.
In SP-RSD, 29.0% of submitted rank-order lists contain errors.
The most common error 41For each group, we take the mean difference between the second-highest bid and the second-highest value. This produces one observation per group playing 2P+X, for a total of 18 observations, and we use a t-test for the null that these have zero mean.
42p = .173 for AC, p = .694 for 2P+X, and p = .290 for AC+X.
26 Figure 3: Histogram: 2nd-highest bid minus 2nd-highest value for AC+X and 2P+X Table 2: mean(abs(2nd bid - 2nd value)) Format Rounds SP OSP p-value Auction 1-5 8.04 3.19 .006 (1.25) (1.05) 6-10 4.99 1.77 .016 (1.18) (0.33) +X Auction 1-5 3.99 1.83 .006 (0.60) (0.41) 6-10 3.69 1.29 .017 (0.87) (0.33) For each group, we take the mean absolute difference over each 5-round block. We then compute standard errors counting each group’s 5-round mean as a single observation. (18 observations per cell, standard errors in parentheses.) p-values are computed using a two-sample t-test, allowing for unequal variances. Other empirical strategies yield similar results; see Appendix B for details.
27 Table 3: Proportion of serial dictatorships not ending in dominant strategy outcome SP OSP p-value Rounds 1-5 43.3% 7.8% .0002 (7.3%) (3.3%) Rounds 6-10 28.9% 6.7% .0011 (5.2%) (3.2%) p-value .1026 .7492 For each group, for each 5-round block, we record the error rate. We then compute standard errors counting each group’s observed error rate as a single observation. (18 observations per cell, standard errors in parentheses.) When comparing SP to OSP, we compute p-values using a two-sample t-test, allowing for unequal variances. (Alternative empirical strategies yield similar results. See Appendix B for details.) When comparing early to late rounds of the same game, we compute p-values using a paired t-test.
under SP-RSD is to swap the ranks of the highest and second-highest prizes, and report the list in order 2nd-1st-3rd-4th. This accounts for 38 out of 209 incorrect rank-order lists. However, errors are diverse: No permutation of {1st, 2nd, 3rd, 4th} accounts for more than a fifth of the incorrect rank-order lists.
In summary, subjects play the dominant strategy at higher rates in OSP mechanisms, as com-pared to SP mechanisms that should (according to standard theory) implement the same allocation rule. This difference is significant and substantial across all three pairs of mechanisms, and persists after five rounds of repeated play with feedback. Of course, the set of all strategy-proof mechanisms cannot be tested exhaustively in a single experiment. This is only preliminary evidence that obvious dominance correctly classifies strategy-proof mechanisms according to their cognitive complexity.
There are not (yet) well-formed alternative theories that explain the data presented here. The following theories predict dominant-strategy play in any strategy-proof mechanism: Level-k equi-librium43, cognitive hierarchy equilibrium44, cursed equilibrium45, and analogy-based expectation equilibrium46. Quantal response equilibrium predicts systematic under-bidding in AC and AC+X, which does not appear in the data.47 6 Discussion In this paper, I proposed a definition of obviously strategy-proof mechanisms. One characterization theorem states: A strategy is obviously dominant if and only if a cognitively limited agent can recognize it as weakly dominant. Another characterization theorem states: A choice rule is OSP-implementable if and only if it can be supported by bilateral commitments. For binary allocation problems, I characterized the OSP mechanisms and the OSP-implementable allocation rules. I produced an impossibility result for a classic matching algorithm.
Games can be simple or complex in many ways, and no formal definition can capture all of them. In particular, since obvious dominance invokes only standard game-theoretic primitives, it does not capture framing effects. The following game has an obviously dominant strategy: You may choose any 100-digit string. You will receive $10 if your string consists only of the digit “7”, and $0 otherwise. This modified game also has an obviously dominant strategy: You may choose 43Stahl and Wilson (1994, 1995); Nagel (1995); Camerer et al. (2004); Crawford and Iriberri (2007a,b).
44Camerer et al. (2004).
45Eyster and Rabin (2005); Esponda (2008).
46Jehiel (2005).
47McKelvey and Palfrey (1995, 1998). See Appendix C for details.
28 any 100-digit string. You will receive $10 if your string consists of the first 100 digits of π, and $0 otherwise. These two games have the same extensive form, so the theory of obvious dominance cannot distinguish them.
Nonetheless, a formal standard of cognitive simplicity is valuable for several reasons. Firstly, a formal standard helps us to make simplicity an explicit design goal, by asking, “What is the optimal simple mechanism for this setting?” Secondly, a formal standard allows us to quantify trade-offs between simplicity and other design goals. For instance, one justification for using a complex mechanism is that no simple mechanism performs well for the problem at hand. Thirdly, a formal standard aids mutual understanding, since our definition of simplicity can be common knowledge, rather than relying on disparate individual intuitions.
Mechanism design typically assumes that the planner can make binding and credible promises, even about events that the promisee does not observe. Sometimes the full commitment assumption is justified, and the literature contains many excellent results for that case. However, sometimes the full commitment assumption is not justified, and we must make do with only partial commitment power. By studying OSP-implementation, we discover which standard results in mechanism design rely sensitively on the assumption of full commitment power, and learn how to design mechanisms that rely only on bilateral commitments.48 Much remains to be done. There are many classic results for SP-implementation, where OSP-implementation is an open question. For instance, in combinatorial auctions, the Vickrey-Clarke-Groves mechanism delivers first-best expected welfare but is not obviously strategy-proof (Vickrey, 1961; Clarke, 1971; Groves, 1973).
However, when objects are substitutes, there always exists an obviously strategy-proof mechanism that delivers at least half of first-best expected welfare (Feldman et al., 2014).49 A natural open question is: What is the welfare-maximizing obviously strategy-proof mechanism for combinatorial auctions?
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32 A Proofs omitted from the main text A.1 Theorem 1 Proof. First we prove the “if” direction. Fix agent 1 and preferences θ1. Suppose that S1 is not obviously dominant in G = ⟨H, ≺, A, A, P, δc, (Ii)i∈N, g⟩. We need to demonstrate that there exists ˜ G that is i-indistinguishable from G, such that λG, ˜ G(S1) is not weakly dominant in ˜ G . We proceed by construction.
Let (S′ 1, I1, hsup, Ssup −1 , dsup c , hinf, Sinf −1, dinf c ) be such that I1 ∈α(S1, S′ 1), hinf ∈I1, hsup ∈I1, and uG 1 (hsup, S′ 1, Ssup −1 , dsup c , θ1) > uG 1 (hinf, S1, Sinf −1, dinf c , θ1) (30) Since G is a game of perfect recall, we can pick (Sinf −1, dinf c ) such that hinf ≺zG(h∅, S1, Sinf −1, dinf c ), by specifying that (Sinf −1, dinf c ) plays in a way consistent with hinf at any h ≺hinf. Likewise for hsup and (Ssup −1 , dsup c ). Suppose we have so done.
We now define another game ˜ G is 1-indistinguishable from G. Intuitively, we construct this as follows: 1. We add a chance move at the start of the game; chance can play L or R.
2. Agent 1 does not at any history know whether chance played L or R.
3. If chance plays L, then the game proceeds as in G.
4. If chance plays R, then the game proceeds mechanically as though all players in N \ 1 and chance played according to Ssup −1 , dsup c in G, with one exception: 5. If chance played R, we reach the information set corresponding to I1, and agent 1 plays S1(I1), then the game henceforth proceeds mechanically as though all players in N \ 1 and chance played according to Sinf −1, dinf c in G.
Formally, the construction proceeds as such: ˜ A = A ∪{L, R}, where A ∩{L, R} = ∅. There is a new starting history ˜ h∅, with two successors σ(˜ h∅) = {˜ hL, ˜ hR}, ˜ A(˜ hL) = L, ˜ A(˜ hR) = R, ˜ P(˜ h∅) = c. The subtree ˜ HL ⊂˜ H starting from ˜ hL ordered by ˜ ≺is the same as the arborescence (H, ≺). ( ˜ A, ˜ P, ˜ δc, ˜ g) are defined on ˜ HL exactly as (A, P, δc, g) are on H. For j ̸= 1, ˜ Ij is defined as on H.
We now construct the subtree starting from ˜ hR.
Let h∗be such that h∗∈σ(hsup), h∗⪯ zG(hsup, S1, Ssup −1 , dsup c ).
H′ ≡{h ∈H : ∃S′′ 1 : h ⪯zG(h∅, S′′ 1, Ssup −1 , dsup c )}} ∩[{h ∈H : P(h) = 1} ∪{h ∈Z}] {h ∈H : h∗⪯h} (31) In words, these are the histories that can be reached by some S′′ 1 when facing Ssup −1 , dsup c , where either agent 1 is called to play or that history is terminal, and such that those histories are not h∗ or its successors.
Let h∗∗be such that h∗∗∈σ(hinf), h∗∗⪯zG(hinf, S1, Sinf −1, dinf c ).
33 H′′ ≡{h ∈H : ∃S′′ 1 : h ⪯zG(h∅, S′′ 1, Sinf −1, dinf c )} ∩[{h ∈H : P(h) = 1} ∪{h ∈Z}] ∩{h ∈H : h∗∗⪯h} (32) In words, these are the histories that can be reached by some S′′ 1 when facing Sinf −1, dinf c , where either agent 1 is called to play or that history is terminal, and such that those histories are h∗∗or its successors.
We now paste these together. Let ˜ HR be the rooted subtree ordered by ˜ ≺, for some bijection γ : ˜ HR →H′ ∪H′′, such that for all ˜ h, ˜ h′ ∈˜ HR, ˜ h˜ ≺˜ h′ if and only if 1. EITHER: γ(˜ h), γ(˜ h′) ∈H′ and γ(˜ h) ≺γ(˜ h′) 2. OR: γ(˜ h), γ(˜ h′) ∈H′′ and γ(˜ h) ≺γ(˜ h′) 3. OR: γ(˜ h) ≺h∗and h∗∗⪯γ(˜ h′) The root of this subtree exists and is unique; it corresponds to γ−1(h), where h is the earliest history preceding zG(h∅, S1, Ssup −1 , dsup c )} where 1 is called to play. Let ˜ hR be the root of ˜ HR. This completes the specification of ˜ H.
For all ˜ h ∈˜ HR, we define: 1. ˜ g(˜ h) = g(γ(˜ h)) if ˜ h is a terminal history.
2. ˜ P(˜ h) = 1 if ˜ h is not a terminal history.
For all ˜ h ∈˜ HR \ ˜ hR, we define ˜ A(˜ h) = A(h), for the unique (˜ h′, h) such that: 1. ˜ h ∈σ(˜ h′) 2. h ∈σ(γ(˜ h′)) 3. h ⪯γ(˜ h) We now specify the information sets for agent 1. Every ˜ h ∈˜ HL corresponds to a unique history in H. We use γL to denote the bijection from ˜ HL to H. Let ˆ γ be defined as γL on ˜ HL and γ on ˜ HR.
1’s information partition ˜ I1 is defined as such: ∀˜ h, ˜ h′ ∈˜ H : ∃˜ I′ 1 ∈˜ I1 : ˜ h, ˜ h′ ∈˜ I′ 1 (33) if and only if ∃I′ 1 ∈I1 : ˆ γ(˜ h), ˆ γ(˜ h′) ∈˜ I′ 1 (34) All that remains is to define δc; we need only specify that at ˜ h∅, c plays R with certainty.50 ˜ G = ⟨˜ H, ˜ ≺, ˜ A, ˜ A, ˜ P, ˜ δc, (˜ Ii)i∈N, ˜ g⟩is 1-indistinguishable from G.
Every experience at some history in ˜ HL corresponds to some experience in G, and vice versa. Moreover, any experience at some history in ˜ HR could also be produced by some history in ˜ HL.
50If one prefers to avoid δc without full support, an alternative proof puts ϵ probability on L, for ϵ close to 0, or for games with |N| > 2 assigns ˜ P(˜ h∅) = 2.
34 Let λG, ˜ G be the appropriate bijection from 1’s information sets and actions in G onto 1’s information sets and actions in ˜ G. Take arbitrary ˜ S−1. Observe that since I1 ∈α(S1, S′ 1), λG, ˜ G(S1) and λG, ˜ G(S′ 1) result in the same histories following ˜ hR, until they reach information set λG, ˜ G(I1).
Having reached that point, λG, ˜ G(S1) leads to outcome g(zG(hinf, S1, Sinf −1, dinf c )) and λG, ˜ G(S′ 1) leads to outcome g(zG(hsup, S′ 1, Ssup −1 , dsup c )). Thus, E˜ δc[u ˜ G 1 (˜ h∅, λG, ˜ G(S′ 1), ˜ S−1, ˜ dc, θ1)] = uG 1 (hsup, S′ 1, Ssup −1 , dsup c , θ1) > uG 1 (hinf, S1, Sinf −1, dinf c , θ1) = E˜ δc[u ˜ G 1 (˜ h∅, λG, ˜ G(S1), ˜ S−1, ˜ dc, θ1)] (35) So λG, ˜ G(S1) is not weakly dominant in ˜ G.
We now prove the “only if” direction. Take arbitrary ˜ G. Suppose λG, ˜ G(S1) ≡˜ S1 is not weakly dominant in ˜ G. We want to show that S1 is not obviously dominant in G.
There exist θ1, ˜ S′ 1 and ˜ S′ −1 such that: E˜ δc[u ˜ G 1 (˜ h∅, ˜ S′ 1, ˜ S′ −1, ˜ dc, θ1)] > E˜ δc[u ˜ G 1 (˜ h∅, ˜ S1, ˜ S′ −1, ˜ dc, θ1)] (36) This inequality must hold for some realization of the chance function, so there exists ˜ dc such that: u ˜ G 1 (˜ h∅, ˜ S′ 1, ˜ S′ −1, ˜ dc, θ1) > u ˜ G 1 (˜ h∅, ˜ S1, ˜ S′ −1, ˜ dc, θ1) (37) Fix ( ˜ S′ 1, ˜ S′ −1, ˜ dc, θ1).
z ˜ G(˜ h∅, ˜ S′ 1, ˜ S′ −1, ˜ dc) ̸= z ˜ G(˜ h∅, ˜ S1, ˜ S′ −1, ˜ dc) (38) Define: ˜ H∗≡{˜ h ∈˜ H : ˜ h ≺z ˜ G(˜ h∅, ˜ S′ 1, ˜ S′ −1, ˜ dc) and ˜ h ≺z ˜ G(˜ h∅, ˜ S1, ˜ S′ −1, ˜ dc)} (39) ˜ h∗≡˜ h ∈˜ H∗: ∀˜ h′ ∈˜ H∗: ˜ h′ ⪯˜ h (40) Since the opponent strategies and chance moves are held constant across both sides of Equation 38, P(˜ h∗) = 1 and ˜ h∗∈˜ I1, where ˜ S1(˜ I1) ̸= ˜ S′ 1(˜ I1). Moreover, ˜ I1 ∈α( ˜ S1, ˜ S′ 1) and λ ˜ G,G(˜ I1) ∈ α(S1, S′ 1), where we denote S′ 1 ≡λ ˜ G,G( ˜ S′ 1).
Since G and ˜ G are 1-indistinguishable, consider the experiences λ ˜ G,G(ψ1(z ˜ G(˜ h∅, ˜ S1, ˜ S′ −1, ˜ dc))) and λ ˜ G,G(ψ1(z ˜ G(˜ h∅, ˜ S′ 1, ˜ S′ −1, ˜ dc))).
In G, λ ˜ G,G(ψ1(z ˜ G(˜ h∅, ˜ S1, ˜ S′ −1, ˜ dc))) could lead to outcome ˜ g(z ˜ G(˜ h∅, ˜ S1, ˜ S′ −1, ˜ dc)). We use (Sinf −1, dinf c ) to denote the corresponding opponent strategies and chance realizations that lead to that outcome.
We denote hinf ≡h ∈λ ˜ G,G(˜ I1) : h ≺zG(h∅, S1, Sinf −1, dinf c ) .
In G, λ ˜ G,G(ψ1(z ˜ G(˜ h∅, ˜ S′ 1, ˜ S′ −1, ˜ dc))) could lead to outcome ˜ g(z ˜ G(˜ h∅, ˜ S′ 1, ˜ S′ −1, ˜ dc)). We use (Ssup −1 , dsup c ) to denote the corresponding opponent strategies and chance realizations that lead to that outcome.
We denote hsup ≡h ∈λ ˜ G,G(˜ I1) : h ≺zG(h∅, S′ 1, Ssup −1 , dsup c ) .
35 uG 1 (hsup, S′ 1, Ssup −1 , dsup c , θ1) = u ˜ G 1 (˜ h∅, ˜ S′ 1, ˜ S′ −1, ˜ dc, θ1) > u ˜ G 1 (˜ h∅, ˜ S1, ˜ S′ −1, ˜ dc, θ1) = uG 1 (hinf, S1, Sinf −1, dinf c , θ1) (41) where hsup, hinf ∈λ ˜ G,G(˜ I1) and λ ˜ G,G(˜ I1) ∈α(S1, S′ 1). Thus S1 is not obviously dominant in G.
A.2 Theorem 2 Proof. The key is to see that, for every G ∈G, there is a corresponding ˜ S∆ 0 , and vice versa. We use ¯ S0 to denote the support of ˜ S∆ 0 . In particular, observe the following isomorphism: Table 4: Equivalence between extensive game forms and Planner mixed strategies G ˜ S∆ 0 dc ˜ S0 ∈¯ S0 δc the probability measure specified by ˜ S∆ 0 g(z) for z ∈Z the Planner’s choice of outcome when she ends the game Ii ((mk, Rk, rk)t−1 k=1, mt, Rt) consistent with some ˜ S0 ∈¯ S0 A(Ii) Rt ψi(z) oc i consistent with some ˜ S0 ∈¯ S0 and ˜ SN Information sets in G are equivalent to sequences of past communication ((mk, Rk, rk)t−1 k=1, mt, Rt) under ˜ S∆ 0 . Available actions at some information set A(Ii) are equivalent to acceptable responses Rt. Thus, for any strategy in some game G, we can construct an equivalent strategy given appro-priate ˜ S∆ 0 , and vice versa.
Furthermore, fixing a chance realization dc and agent strategies SN uniquely results in some outcome. Similarly, fixing a realization of the planner’s mixed strategy ˜ S0 ∈¯ S0 and agent strategies ˜ SN uniquely determines some outcome. Consequently, for any G ∈G, there exists ˜ S∆ 0 with the same strategies available for each agent and the same resulting (probability measure over) outcomes, and vice versa.51 Table 4 summarizes.
The next step is to see that a bilateral commitment ˆ Si 0 is equivalent to the Planner promising to ‘run’ only games in some equivalence class that is i-indistinguishable.
Suppose that there is some G that OSP-implements f. Pick some equivalent ˜ S∆ 0 with support ¯ S0.
For each i ∈N, specify the bilateral commitment ˆ Si 0 ≡Φ−1 i (Φi(¯ S0)). These bilateral commitments support f.
To see this, take any ˜ S∆′ 0 ∈∆ˆ Si 0, with support ¯ S′ 0. For any ˜ S′ 0 ∈¯ S′ 0, for any ˜ S′ N, there exists ˜ S0 ∈¯ S0 and ˜ SN such that φi( ˜ S′ 0, ˜ S′ N) = φi( ˜ S0, ˜ SN). By construction, G is such that: There exists z ∈Z where ψi(z) and g(z) are equivalent to φi( ˜ S0, ˜ SN). Thus, for G′ that is equivalent to ˜ S∆′ 0 , every terminal history in G′ results in the same experience for i and the same outcome as some terminal history in G. Consequently, G and G′ are i-indistinguishable. Thus, by Theorem 1, the strategy assigned to agent i with type θi is weakly dominant in G′, which implies that it is a best 51Implicitly, this relies on the requirement that both G and ˜ S∆ 0 have finite length. If one had finite length but the other could be infinitely long, the resulting outcome would not be well defined and the equivalence would not hold.
36 response to ˜ S∆′ 0 and any ˜ SN in the bilateral commitment game. Thus, if f is OSP-implementable, then f can be supported by bilateral commitments.
Suppose that f can be supported by bilateral commitments (ˆ Si 0)i∈N, with requisite ˜ S∆ 0 (with support ¯ S0) and ( ˜ Sθ N)θ∈Θ. Without loss of generality, let us suppose these are ‘minimal’ bilateral commitments, i.e. ˆ Si 0 = Φ−1 i (Φi(¯ S0)). Pick G that is equivalent to ˜ S∆ 0 . G OSP-implements f.
To see this, consider any G′ such that G and G′ are i-indistinguishable. Let ˜ S∆′ 0 denote the Planner strategy that corresponds to G′. At any terminal history z′ in G′, the resulting experience ψi(z′) and outcome g′(z′) are equivalent to the experience ψi(z) and outcome g(z) for some terminal history z in G. These in turn correspond to some observation oi ∈Φi(¯ S0). Thus ˜ S∆′ 0 ∈∆ˆ Si 0. Since f is supported by (ˆ Sj 0)j∈N, ˜ Sθi i is a best response (for type θi) to ˜ S∆′ 0 and any ˜ SN\i. Thus, the equiv-alent strategy Sθi i is weakly dominant in G′. Since this argument holds for all i-indistinguishable G′, by Theorem 1, Sθi i is obviously dominant in G.
Thus, if f can be supported by bilateral commitments, then f is OSP-implementable.
A.3 Proposition 2 Proof. We prove the contrapositive. Suppose ( ˜ G, ( ˜ Sθ)θ∈Θ) does not OSP-implement f. Then there exists some (i, θi, ˜ Sθi i , ˜ S′ i, ˜ Ii) such that ˜ Ii ∈α( ˜ Sθi i , ˜ S′ i) and u ˜ G i (˜ h, ˜ Sθi i , ˜ S−i, ˜ dc, θi, ) < u ˜ G i (˜ h′, ˜ S′ i, ˜ S′ −i, ˜ d′ c, θi, ) (42) for some (˜ h, ˜ S−i, ˜ dc) and (˜ h′, ˜ S′ −i, ˜ d′ c).
Notice that ˜ h and ˜ h′ correspond to histories h and h′ in G. Moreover, we can define S′ i = ˜ S′ i at information sets containing histories that are shared by G and ˜ G, and specify S′ i arbitrarily elsewhere. We do the same for ( ˜ S−i, ˜ dc) and ( ˜ S′ −i, ˜ d′ c), to construct (S−i, dc) and (S−i, dc). But, starting from h and h′ respectively, these result in the same outcomes as their partners in ˜ G. Thus, uG i (h, Sθi i , S−i, dc, θi, ) < uG i (h′, S′ i, S′ −i, d′ c, θi, ) (43) h, h′ ∈Ii, for Ii ∈α(Sθi i , S′ i). Thus (G, (Sθ)θ∈Θ) does not OSP-implement f.
A.4 Theorem 3 Proof. Take any (G, (Sθ)θ∈Θ) that implements (fy, ft). For any history h, we define Θh ≡{θ ∈Θ : h ⪯zG(h∅, Sθ)} (44) Θh,i ≡{θi : ∃θ−i : (θi, θ−i) ∈Θh} (45) For information set Ii, we define ΘIi ≡∪h∈IiΘh (46) ΘIi,i ≡{θi : ∃θ−i : (θi, θ−i) ∈ΘIi} (47) Θ1 Ii,i ≡{θi : ∃θ−i : (θi, θ−i) ∈ΘIi and i ∈fy(θi, θ−i)} (48) 37 Θ0 Ii,i ≡{θi : ∃θ−i : (θi, θ−i) ∈ΘIi and i / ∈fy(θi, θ−i)} (49) Some observations about this construction: 1. Since player i’s strategy depends only on his own type, ΘIi,i = Θh,i for all h ∈Ii.
2. ΘIi,i = Θ1 Ii,i ∪Θ0 Ii,i 3. Since SP requires 1i∈fy(θ) weakly increasing in θi, Θ1 Ii,i dominates Θ0 Ii,i in the strong set order.
Lemma 1. Suppose (G, (Sθ)θ∈Θ) OSP-implements (fy, ft), where G ≡⟨H, ≺, A, A, P, δc, (Ii)i∈N, g⟩.
For all i, for all Ii ∈Ii, if: 1. θi < θ′ i 2. θi ∈Θ1 Ii,i 3. θ′ i ∈Θ0 Ii,i then Sθi i (Ii) = Sθ′ i i (Ii).
Equivalently, for any Ii, there exists a∗ Ii such that for all θi ∈Θ1 Ii,i ∩Θ0 Ii,i, Sθi i (Ii) = a∗ Ii.
Suppose not. Take (i, Ii, θi, θ′ i) constituting a counterexample to Lemma 1. Since θi ∈Θ1 Ii,i, there exists h ∈Ii and S−i such that i ∈gy(zG(h, Sθi i , S−i)). Fix ti ≡gt,i(zG(h, Sθi i , S−i)). Since θ′ i ∈ Θ0 Ii,i, there exists h′ ∈Ii and S′ −i such that i / ∈gy(zG(h′, Sθ′ i i , S′ −i)). Fix t′ i ≡gt,i(zG(h′, Sθ′ i i , S′ −i)).
Since Sθi i (Ii) ̸= Sθ′ i i (Ii) and θi ∪θ′ i ⊆ΘIi,i, Ii ∈α(Sθi i , Sθ′ i i ). Thus, OSP requires that ui(θi, h, Sθi i , S−i) ≥ui(θi, h′, Sθ′ i i , S′ −i) (50) which implies θi + ti ≥t′ i (51) and ui(θ′ i, h, Sθi i , S−i) ≤ui(θ′ i, h′, Sθ′ i i , S′ −i) (52) which implies θ′ i + ti ≤t′ i (53) But θ′ i > θi, so θ′ i + ti > t′ i (54) a contradiction. This proves Lemma 1.
Lemma 2. Suppose (G, (Sθ)θ∈Θ) OSP-implements (fy, ft) and P(G, (Sθ)θ∈Θ) = G. Take any Ii such that Θ1 Ii,i ∩Θ0 Ii,i ̸= ∅, and associated a∗ Ii.
1. If there exists θi ∈Θ0 Ii,i such that Sθi i (Ii) ̸= a∗ Ii, then there exists t0 i such that: 38 (a) For all θi ∈Θ0 Ii,i such that Sθi i (Ii) ̸= a∗ Ii, for all h ∈Ii, for all S−i, gt,i(zG(h, Sθi i , S−i)) = t0 i .
(b) For all θi ∈ΘIi,i such that Sθi i (Ii) = a∗ Ii, for all h ∈Ii, for all S−i, if i / ∈gy(zG(h, Sθi i , S−i)), then gt,i(zG(h, Sθi i , S−i)) = t0 i .
2. If there exists θi ∈Θ1 Ii,i such that Sθi i (Ii) ̸= a∗ Ii, then there exists t1 i such that: (a) For all θi ∈Θ1 Ii,i such that Sθi i (Ii) ̸= a∗ Ii, for all h ∈Ii, for all S−i, gt,i(zG(h, Sθi i , S−i)) = t1 i .
(b) For all θi ∈ΘIi,i such that Sθi i (Ii) = a∗ Ii, for all h ∈Ii, for all S−i, if i ∈gy(zG(h, Sθi i , S−i)), then gt,i(zG(h, Sθi i , S−i)) = t1 i .
Take any type θ′ i ∈Θ0 Ii,i such that Sθ′ i i (Ii) ̸= a∗ Ii. Take any type θ′′ i ∈Θ0 Ii,i such that Sθ′′ i i (Ii) = a∗ Ii.
(By Θ1 Ii,i ∩Θ0 Ii,i ̸= ∅there exists at least one such type.) Notice that Ii ∈α(Sθ′ i i , Sθ′′ i i ).
By Lemma 1, θ′ i / ∈Θ1 Ii,i, and the game is pruned. Thus, ∀h ∈Ii : ∀S−i : i / ∈gy(zG(h, Sθ′ i i , S−i)) (55) Since θ′′ i ∈Θ0 Ii,i, ∃h ∈Ii : ∃S−i : i / ∈gy(zG(h, Sθ′′ i i , S−i)) (56) OSP requires that type θ′ i does not want to (inf-sup) deviate. Thus, inf h∈Ii,S−i gt,i(zG(h, Sθ′ i i , S−i)) ≥ sup h∈Ii,S−i {gt,i(zG(h, Sθ′′ i i , S−i)) : i / ∈gy(zG(h, Sθ′′ i i , S−i))} (57) OSP also requires that type θ′′ i does not want to (inf-sup) deviate. This implies inf h∈Ii,S−i{gt,i(zG(h, Sθ′′ i i , S−i)) : i / ∈gy(zG(h, Sθ′′ i i , S−i))} ≥ sup h∈Ii,S−i gt,i(zG(h, Sθ′ i i , S−i)) (58) The RHS of Equation 57 is weakly greater than the LHS of Equation 58. The RHS of Equation 58 is weakly greater than the LHS of Equation 57. Consequently all four terms are equal. Moreover, this argument applies to every θ′ i ∈Θ0 Ii,i such that Sθ′ i i (Ii) ̸= a∗ Ii, and every θ′′ i ∈Θ0 Ii,i such that Sθ′′ i i (Ii) = a∗ Ii. Since the game is pruned, θ′′ i satisfies (1b) iffθ′′ i ∈Θ0 Ii,i and Sθ′′ i i (Ii) = a∗ Ii. This proves part 1 of Lemma 2. Part 2 follows by symmetry; we omit the details since they involve only small changes to the above argument.
Lemma 3. Suppose (G, (Sθ)θ∈Θ) OSP-implements (fy, ft) and P(G, (Sθ)θ∈Θ) = G. Take any Ii such that Θ1 Ii,i ∩Θ0 Ii,i ̸= ∅, and associated a∗ Ii. Let t1 i and t0 i be defined as before.
1. If there exists θi ∈Θ0 Ii,i such that Sθi i (Ii) ̸= a∗ Ii, then for all (h ∈Ii, Si, S−i), if i ∈ gy(zG(h, Si S−i)), then gt,i(zG(h, Si, S−i)) ≤t0 i −sup{θi ∈Θ0 Ii,i : Sθi i (Ii) ̸= a∗ Ii}.
39 2. If there exists θi ∈Θ1 Ii,i such that Sθi i (Ii) ̸= a∗ Ii, then for all (h ∈Ii, Si, S−i), if i / ∈ gy(zG(h, Si S−i)), then gt,i(zG(h, Si, S−i)) ≤inf{θi ∈Θ1 Ii,i : Sθi i (Ii) ̸= a∗ Ii} + t1 i .
Suppose that part 1 of Lemma 3 does not hold. Fix (h ∈Ii, Si, S−i) such that i ∈gy(zG(h, Si S−i)) and gt,i(zG(h, Si, S−i)) > t0 i −sup{θi ∈Θ0 Ii,i : Sθi i (Ii) ̸= a∗ Ii}. Since G is pruned, we can find some θ′ i ∈ΘIi,i such that for every ˜ Ii ∈{I′ i ∈Ii : Ii ∈ψ(I′ i)}, Sθ′ i i (˜ Ii) = Si(˜ Ii). Fix that θ′ i.
Fix θ′′ i ∈Θ0 Ii,i such that Sθi i (Ii) ̸= a∗ Ii and θ′′ i ≥sup{θi ∈Θ0 Ii,i : Sθi i (Ii) ̸= a∗ Ii} −ϵ. Since G is pruned and θ′′ i / ∈Θ1 Ii,i (by Lemma 1), it must be that Sθ′′ i i (Ii) ̸= Sθ′ i i (Ii).
By construction, Ii ∈α(Sθ′ i i , Sθ′′ i i ).
OSP requires that, for all h′′ ∈Ii, S′′ −i: ui(θ′′ i , h′′, Sθ′′ i i , S′′ −i) ≥ui(θ′′ i , h, Sθ′ i i , S−i) (59) which entails t0 i ≥θ′′ i + gt,i(zG(h, Si, S−i)) (60) which entails t0 i −sup{θi ∈Θ0 Ii,i : Sθi i (Ii) ̸= a∗ Ii} + ϵ ≥gt,i(zG(h, Si, S−i)) (61) But, by hypothesis, t0 i −sup{θi ∈Θ0 Ii,i : Sθi i (Ii) ̸= a∗ Ii} < gt,i(zG(h, Si, S−i)) (62) Since this argument holds for all ϵ > 0, we can pick ϵ small enough to create a contradiction.
This proves part 1 of Lemma 3. Part 2 follows by symmetry.
Lemma 4. Suppose (G, (Sθ)θ∈Θ) OSP-implements (fy, ft) and P(G, (Sθ)θ∈Θ) = G. Take any Ii such that |Θ1 Ii,i ∩Θ0 Ii,i| > 1 and associated a∗ Ii.
1. If there exists θi ∈Θ0 Ii,i such that Sθi i (Ii) ̸= a∗ Ii, then for all θ′ i ∈Θ1 Ii,i, Sθ′ i i (Ii) = a∗ Ii.
2. (Equivalently) If there exists θi ∈Θ1 Ii,i such that Sθi i (Ii) ̸= a∗ Ii, then for all θ′ i ∈Θ0 Ii,i, Sθ′ i i (Ii) = a∗ Ii.
Suppose Part 1 of Lemma 4 does not hold. Fix Ii, and choose θ′ i < θ′′ i such that {θ′ i} ∪{θ′′ i } ⊆ Θ1 Ii,i ∩Θ0 Ii,i. Fix θ′′′ i ∈Θ1 Ii,i such that Sθ′′′ i i (Ii) ̸= a∗ Ii. By Lemma 1, if θ′′′ i ∈Θ0 Ii,i, then Sθ′′′ i i (Ii) = a∗ Ii, a contradiction. Thus, θ′′′ i ∈Θ1 Ii,i \ Θ0 Ii,i, and since Θ1 Ii,i dominates Θ0 Ii,i in the strong set order, θ′′ i < θ′′′ i .
Since θ′ i ∈Θ1 Ii,i, there exists h′ ∈Ii and θ′ −i such that (θ′ i, θ′ −i) ∈ΘIi and i ∈gy(zG(h′, Sθ′ i i , S θ′ −i −i )).
By Lemma 2, there exists ai ∈Ai(Ii) such that ai ̸= a∗ Ii and choosing ai ensures i / ∈y and ti = t0 i .
Thus, by G SP θ′ i + gt,i(zG(h′, Sθ′ i i , S θ′ −i −i )) ≥t0 i (63) By θ′′ i ∈Θ0 Ii,i, there exists h′′ ∈Ii and θ′′ −i such that i / ∈gy(zG(h′′, Sθ′′ i i , S θ′′ −i −i )). By Lemma 2 gt,i(zG(h′′, Sθ′′ i i , S θ′′ −i −i )) = t0 i (64) 40 By G SP, i ∈gy(zG(h′, Sθ′′′ i i , S θ′ −i −i )) and gt,i(zG(h′, Sθ′′′ i i , S θ′ −i −i )) = gt,i(zG(h′, Sθ′ i i , S θ′ −i −i )).
Notice that Ii ∈α(Sθ′′ i i , Sθ′′′ i i ). Thus, OSP requires that θ′′ i does not want to (inf-sup) deviate to θ′′′ i ’s strategy, which entails: gt,i(zG(h′′, Sθ′′ i i , S θ′′ −i −i )) ≥θ′′ i + gt,i(zG(h′, Sθ′′′ i i , S θ′ −i −i )) (65) t0 i ≥θ′′ i + gt,i(zG(h′, Sθ′′′ i i , S θ′ −i −i )) > θ′ i + gt,i(zG(h′, Sθ′ i i , S θ′ −i −i )) (66) which contradicts Equation 63.
Part 2 is the contrapositive of Part 1. This proves Lemma 4.
Lemma 5. Suppose (G, (Sθ)θ∈Θ) OSP-implements (fy, ft) and P(G, (Sθ)θ∈Θ) = G.
For all Ii, if |Θ1 Ii,i ∩Θ0 Ii,i| ≤1 and |A(Ii)| ≥2, then there exists t1 i and t0 i such that: 1. For all θi ∈ΘIi,i, h ∈Ii, S−i: (a) If i / ∈gy(zG(h, Sθi i , S−i)) then gt,i(zG(h, Sθi i , S−i)) = t0 i (b) If i ∈gy(zG(h, Sθi i , S−i)) then gt,i(zG(h, Sθi i , S−i)) = t1 i 2. If |Θ1 Ii,i| > 0 and |Θ0 Ii,i| > 0, then t1 i = −inf{θi ∈Θ1 Ii,i} + t0 i By G pruned, ΘIi,i ̸= ∅.
Consider the case where Θ1 Ii,i = ∅.
Pick some θ′ i ∈Θ0 Ii,i and some h′ ∈Ii, S′ −i.
Fix t0 i ≡ gt,i(zG(h′, Sθ′ i i , S′ −i)). Suppose there exists some (θ′′ i , θ′′ −i) ∈ΘIi such that ft,i(θ′′ i , θ′′ −i) = t0 i ′′ ̸= t0 i .
Pick h′′ ∈Ii such that h′′ ⪯zG(∅, Sθ′′ i i , S θ′′ −i −i ). By Equation 12, for all θi ∈ΘIi,i, ft,i(θi, θ′′ −i) = t0 i ′′.
By G pruned and |A(Ii)| ≥2, we can pick θ′′′ i ∈Θ0 Ii,i such that Sθ′′′ i i (Ii) ̸= Sθ′ i i (Ii). Notice that Ii ∈α(Sθ′′′ i i , Sθ′ i i ). If t0 i ′′ > t0 i , then ui(θ′ i, h′, Sθ′ i i , S′ −i) = t0 i < t0 i ′′ = ui(θ′ i, h′′, Sθ′′′ i i , S θ′′ −i −i ) (67) so Sθ′ i i is not obviously dominant for (i, θ′ i). If t0 i ′′ < t0 i , then ui(θ′′′ i , h′, Sθ′ i i , S′ −i) = t0 i > t0 i ′′ = ui(θ′′′ i , h′′, Sθ′′′ i i , S θ′′ −i −i ) (68) so Sθ′′′ i i is not obviously dominant for (i, θ′′′ i ). By contradiction, this proves Lemma 5 for this case. A symmetric argument proves Lemma 5 for the case where Θ0 Ii,i = ∅.
Note that, if Lemma 5 holds at some information set Ii, it holds at all information sets I′ i that follow Ii. Thus, we need only consider some earliest information set I∗ i at which |Θ1 I∗ i ,i ∩Θ0 I∗ i ,i| ≤1 and |A(I∗ i )| ≥2.
Now we consider the case where Θ1 I∗ i ,i ̸= ∅and Θ0 I∗ i ,i ̸= ∅.
At every prior information set Ii prior to I∗ i , |Θ1 Ii,i ∩Θ0 Ii,i| > 1. Since Θ1 I∗ i ,i ̸= ∅and Θ0 I∗ i ,i ̸= ∅, by Lemma 4, I∗ i is reached by some interval of types all taking the same action. Thus sup{θi ∈ Θ0 I∗ i ,i} = inf{θi ∈Θ1 I∗ i ,i}.
41 Fix ˆ θi ∈Θ0 I∗ i ,i such that ˆ θi ≥sup{θi ∈Θ0 I∗ i ,i}−ϵ. Choose corresponding ˆ h ∈I∗ i and ˆ θ−i ∈ΘI∗ i ,−i such that i / ∈gy(zG(ˆ h, S ˆ θi i , S ˆ θ−i −i )). Define t0 i ≡gt,i(zG(ˆ h, S ˆ θi i , S ˆ θ−i −i )).
Fix ˇ θi ∈Θ1 I∗ i ,i such that ˇ θi ≤inf{θi ∈Θ1 I∗ i ,i} + ϵ. Choose corresponding ˇ h ∈I∗ i and ˇ θ−i ∈ΘI∗ i ,−i such that i ∈gy(zG(ˇ h, S ˇ θi i , S ˇ θ−i −i )). Define t1 i ≡gt,i(zG(ˇ h, S ˇ θi i , S ˇ θ−i −i )).
Suppose there exists some (θ′ i, θ′ −i) ∈ΘI∗ i such that i / ∈fy(θ′ i, θ′ −i) and ft,i(θ′ i, θ′ −i) = t0 i ′ ̸= t0 i .
Since sup{θi ∈Θ0 I∗ i ,i} = inf{θi ∈Θ1 I∗ i ,i}, it follows that for all θ−i ∈ΘI∗ i ,−i, inf{θi : i ∈fy(θi, θ−i)} = inf{θi ∈Θ1 I∗ i ,i}.
Thus, by Equation 12, for all θi ∈ΘI∗ i ,i: ft,i(θi, θ′ −i) = −1i∈fy(θi,θ′ −i) inf{θi ∈ Θ1 I∗ i ,i} + t0 i ′. Fix h′ ∈I∗ i such that h′ ⪯zG(∅, Sθ′ i i , S θ′ −i −i ).
By A(Ii) ≥2, we can pick some θ′′ i ∈ΘI∗ i ,i such that Sθ′′ i i (Ii) ̸= S ˆ θi i (Ii). Notice that I∗ i ∈ α(Sθ′′ i i , S ˆ θi i ). Either θ′′ i ∈Θ0 I∗ i or θ′′ i ∈Θ1 I∗ i \ Θ0 I∗ i . Suppose θ′′ i ∈Θ0 I∗ i . Suppose t0 i ′ > t0 i . By OSP, ui(ˆ θi, ˆ h, S ˆ θi i , S ˆ θ−i −i ) ≥ui(ˆ θi, h′, Sθ′′ i i , S θ′ −i −i ) (69) which entails t0 i ≥t0 i ′ + 1i∈fy(θ′′ i ,θ′ −i)(ˆ θi −inf{θi ∈Θ1 I∗ i ,i}) ≥t0 i ′ + 1i∈fy(θ′′ i ,θ′ −i)(−ϵ) ≥t0 i ′ −ϵ (70) and we can pick ϵ small enough to constitute a contradiction. Suppose t0 i ′ < t0 i . By OSP ui(θ′′ i , ˆ h, S ˆ θi i , S ˆ θ−i −i ) ≤ui(θ′′ i , h′, Sθ′′ i i , S θ′ −i −i ) (71) which entails t0 i ≤t0 i ′ + 1i∈fy(θ′′ i ,θ′ −i)(θ′′ i −inf{θi ∈Θ1 I∗ i ,i}) = t0 i ′ + 1i∈fy(θ′′ i ,θ′ −i)(θ′′ i −sup{θi ∈Θ0 I∗ i ,i}) ≤t0 i ′ (72) which is a contradiction.
The case that remains is θ′′ i ∈Θ1 I∗ i \ Θ0 I∗ i . Then i ∈fy(θ′′ i , θ′ −i) and ft,i(θ′′ i , θ′ −i) = −inf{θi ∈ Θ1 I∗ i ,i} + t0 i ′. Suppose t0 i ′ > t0 i . OSP requires: ui(ˆ θi, ˆ h, S ˆ θi i , S ˆ θ−i −i ) ≥ui(ˆ θi, h′, Sθ′′ i i , S θ′ −i −i ) (73) which entails t0 i ≥ˆ θi −inf{θi ∈Θ1 I∗ i ,i} + t0 i ′ ≥t0 i ′ −ϵ (74) and we can pick ϵ small enough to constitute a contradiction.
Suppose t0 i ′ < t0 i . Since Sθ′′ i i (Ii) ̸= S ˆ θi i (Ii), either S ˇ θi i (Ii) ̸= S ˆ θi i (Ii) or S ˇ θi i (Ii) ̸= Sθ′′ i i (Ii). More-over, ft,i(ˇ θi, θ′ i) = −1i∈fy(ˇ θi,θ′ −i) inf{θi ∈Θ1 I∗ i ,i} + t0 i ′.
42 Suppose S ˇ θi i (Ii) ̸= S ˆ θi i (Ii). OSP requires: ui(ˇ θi, h′, S ˇ θi i , S θ′ −i −i ) ≥ui(ˇ θi, ˆ h, S ˆ θi i , S ˆ θ−i −i ) (75) which entails 1i∈fy(ˇ θi,θ′ −i)(ˇ θi −inf{θi ∈Θ1 I∗ i ,i}) + t0 i ′ ≥t0 i (76) which entails 1i∈fy(ˇ θi,θ′ −i)ϵ + t0 i ′ ≥t0 i (77) and we can pick ϵ small enough to yield a contradiction. Suppose S ˇ θi i (Ii) ̸= Sθ′′ i i (Ii). By Equation 12, ft,i(ˇ θi, θ′ −i) = −1i∈fy(ˇ θi,θ′ −i) inf{θi ∈Θ1 I∗ i ,i} + t0 i ′, and ft,i(θ′′ i , ˆ θ−i) = −inf{θi ∈Θ1 I∗ i ,i} + t0 i . OSP requires ui(ˇ θi, h′, S ˇ θi i , S θ′ −i −i ) ≥ui(ˇ θi, ˆ h, Sθ′′ i i , S ˆ θ−i −i ) (78) which entails 1i∈fy(ˇ θi,θ′ −i)(ˇ θi −inf{θi ∈Θ1 I∗ i ,i}) + t0 i ′ ≥(ˇ θi −inf{θi ∈Θ1 I∗ i ,i}) + t0 i (79) which entails 1i∈fy(ˇ θi,θ′ −i)ϵ + t0 i ′ ≥ϵ + t0 i (80) which entails t0 i ′ ≥t0 i (81) a contradiction. By the above argument, for all Ii satisfying the assumptions of Lemma 5, there is a unique transfer t0 i for all terminal histories z passing through Ii such that i / ∈gy(z). Equation 12 thus implies that there is a unique transfer t1 i for all terminal histories z passing through Ii such that i ∈gy(z). Moreover, t1 i = −inf{θi ∈Θ1 Ii,i} + t0 i . This proves Lemma 5.
Now to bring this all together. We leave showing parts (1.c.iv) and (2.c.iv) of Definition 18 to the last. Take any (G, (Sθ)θ∈Θ) that OSP-implements (fy, ft). Define ˜ G ≡P(G, (Sθ)θ∈Θ) and ( ˜ Sθ)θ∈Θ as (Sθ)θ∈Θ restricted to ˜ G. By Proposition 2, ( ˜ G, ( ˜ Sθ)θ∈Θ) OSP-implements (fy, ft).
We now characterize ( ˜ G, ( ˜ Sθ)θ∈Θ). For any player i, consider any information set I∗ i such that |A(I∗ i )| ≥2, and for all prior information sets I′ i ≺I∗ i , |A(I′ i)| = 1. By Lemma 1, there is a unique action a∗ I∗ i taken by all types in Θ1 I∗ i ,i ∩Θ0 I∗ i ,i.
Either |Θ1 I∗ i ,i ∩Θ0 I∗ i ,i| > 1 or |Θ1 I∗ i ,i ∩Θ0 I∗ i ,i| ≤1.
If |Θ1 I∗ i ,i ∩Θ0 I∗ i ,i| > 1, then by Lemma 4, ˜ G pruned and |A(Ii)| ≥2, 1. EITHER: There exists θi ∈Θ0 I∗ i ,i such that Sθi i (I∗ i ) ̸= a∗ I∗ i , and for all θ′ i ∈Θ1 I∗ i ,i, Sθ′ i i (I∗ i ) = a∗ I∗ i .
2. OR: There exists θi ∈Θ1 I∗ i ,i such that Sθi i (I∗ i ) ̸= a∗ I∗ i , and for all θ′ i ∈Θ0 I∗ i ,i, Sθ′ i i (I∗ i ) = a∗ I∗ i .
43 In the first case, then by Lemma 2, there is some t0 i such that, for all (Si, S−i), for all h ∈Ii, if i / ∈gy(zG(h, Si, S−i), then gt,i(zG(h, Si, S−i) = t0 i . Moreover, we can define a ‘going transfer’ at all information sets I′ i such that I∗ i ⪯I′ i: ˜ t1 i (I′ i) ≡ min I′′ i :I∗ i ⪯I′′ i ⪯I′ i [t0 i −sup{θi ∈Θ0 I′′ i ,i : Sθi i (I′′ i ) ̸= a∗ I′′ i }] (82) Notice that this function falls monotonically as we move along the game tree; for any I′ i, I′′ i such that I′ i ⪯I′′ i , ˜ t1 i (I′ i) ≥˜ t1 i (I′′ i ). Moreover, by construction, at any I′ i, I′′ i such that I′′ i is the immediate successor of I′ i in i’s experience, if ˜ t1 i (I′ i) > ˜ t1 i (I′′ i ), then there exists a ∈A(I′′ i ) that yields i / ∈y, and by Lemma 2 this yields transfer t0 i . We define A0 to include all such quitting actions; i.e. A0 is the set of all actions such that: 1. a ∈Ii for some Ii ∈Ii 2. For all z such that a ∈ψi(z): i / ∈gy(z) and gt,i(z) = t0 i Lemma 3 and SP together imply that, at any terminal history z, if i ∈gy(z), then gt,i(z) = inf Ii:I∗ i ⪯Ii≺z ˜ t1 i (Ii) (83) This holds because, if gt,i(z) < infIi:I∗ i ⪯Ii≺z ˜ t1 i (Ii), then type θi such that t0 i −infIi:I∗ i ⪯Ii≺z ˜ t1 i (Ii)) < θi < t0 i −gt,i(z) could profitably deviate to play a ∈A0 at information set I∗ i .
In the second case, then by Lemma 2, there is some t1 i such that, for all (Si, S−i), for all h ∈I∗ i , if i ∈gy(zG(h, Si, S−i), then gt,i(zG(h, Si, S−i) = t1 i . Moreover, we can define a ‘going transfer’ at all information sets I′ i such that I∗ i ⪯I′ i): ˜ t0 i (I′ i) ≡ min I′′ i :I∗ i ⪯I′′ i ⪯I′ i [t1 i + inf{θi ∈Θ0 I′′ i ,i : Sθi i (I′′ i ) ̸= a∗ I′′ i }] (84) Once more, 1. This function falls monotonically as we move along the game tree.
2. At any I′ i, I′′ i such that I′′ i is the immediate successor of I′ i in i’s experience, if ˜ t0 i (I′ i) > ˜ t0 i (I′′ i ), then there exists a ∈A(I′′ i ) that yields i ∈y, and transfer t1 i .
3. For any z, if i / ∈gy(z), then gt,i(z) = inf Ii:I∗ i ⪯Ii≺z ˜ t0 i (Ii) (85) We define A1 symmetrically for this second case.
Part (1.c.iii) and (2.c.iii) of Definition 18 follow from Lemma 5. The above constructions suffice to prove Theorem 3 for cases where |Θ1 I∗ i ,i ∩Θ0 I∗ i ,i| > 1. Cases where |Θ1 I∗ i ,i ∩Θ0 I∗ i ,i| ≤1 are dealt with by Lemma 5.
Now for the last piece: We prove that parts (1.c.iv) and (2.c.iv) of Definition 18 hold. The proof of part (1.c.iv) is as follows: Suppose we are facing the “either” clause of Definition 18, and for some Ii, |A(I′ i) \ A0| > 1. By part (1.c.iii), we know that the going transfer ˜ t1 i can fall no further.
Since ˜ G is pruned and |A(I′ i) \ A0| > 1, there exist two distinct types of i, θi, θ′ i ∈ΘI′ i,i, who do not quit at I′ i, and take different actions. Since neither quits at I′ i and the going transfer falls no further, there exist θ−i, θ′ −i ∈ΘI′ i,−i such that i ∈fy(θi, θ−i) and i ∈fy(θ′ i, θ′ −i). So there exist (h ∈I′ i, S−i) and (h′ ∈I′ i, S′ −i) such that 44 i ∈gy(h, Sθi i , S−i) (86) i ∈gy(h′, Sθ′ i i , S′ −i) (87) gt,i(h, Sθi i , S−i) = gt,i(h′, Sθ′ i i , S′ −i) = ˜ t1 i (I′ i) (88) WLOG suppose θi < θ′ i. Suppose that there does not exist a ∈A(I′ i) such that, for all z such that a ∈ψi(z), i ∈gy(z). Then there must exist (h′′ ∈I′ i, S′′ −i) such that i / ∈gy(h′′, Sθ′ i i , S′′ −i) (89) gt,i(h′′, Sθ′ i i , S′′ −i) = t0 i (90) Note that I′ i ∈α(Sθi i , Sθ′ i i ). But then Sθ′ i i is not obviously dominant, a contradiction, since u ˜ G i (h′′, Sθ′ i i , S′′ −i, θ′ i) = t0 i ≤θi + ˜ t1 i (I′ i) < θ′ i + ˜ t1 i (I′ i) = u ˜ G i (h, Sθi i , S−i, θ′ i) (91) (The first inequality holds because of type θi’s incentive constraint.) This shows that part (1.c.iv) of Definition 18 holds. Part (2.c.iv) is proved symmetrically.
A.5 Theorem 4 Proof. Take any monotone price mechanism G. For any i, the following strategy Si is obviously dominant: 1. If i encounters an information set consistent with Clause 1 of Definition 18, then, from that point forward: (a) If θi + ˜ t1 i (Ii) > t0 i and there exists a ∈A(Ii) \ A0, play a ∈A(Ii) \ A0.
i. If |A(Ii) \ A0| > 1, then play a ∈A(I′ i) such that: For all z such that a ∈ψi(z): i ∈gy(z).
(b) Else play some a ∈A0.52 2. If i encounters an information set consistent with Clause 2 of Definition 18, then, from that point forward: (a) If θi + t1 i < ˜ t0 i (Ii) and there exists a ∈A(Ii) \ A1, play a ∈A(Ii) \ A1.
i. If |A(Ii) \ A1| > 1, then play a ∈A(I′ i) such that: For all z such that a ∈ψi(z): i / ∈gy(z).
(b) Else play some a ∈A1.
52If |A0 ∩A(Ii)| > 1, the agent chooses deterministically but arbitrarily.
45 The above strategy is well-defined for any agent in any monotone price mechanism, by inspection of Definition 18.
Consider any deviating strategy S′ i. At any earliest point of departure, the agent will have encountered an information set consistent with either Clause 1 or Clause 2 of Definition 18. Suppose that the agent has encountered an information set covered by Clause 1.
Take some earliest point of departure Ii ∈α(Si, S′ i). Notice that, by (1.d) of Definition 18, no matter what strategy i plays, conditional on reaching Ii, either agent i is not in the allocation and receives t0 i , or agent i is in the allocation and receives a transfer ˆ ti ≤˜ t1 i (Ii).
Suppose θi + ˜ t1 i (Ii) > t0 i . Note that under Si, conditional on reaching Ii, the agent either is not in the allocation and receives t0 i , or is in the allocation and receives a transfer strictly above t0 i −θi.
If S′ i(Ii) ∈A0 (i.e. if agent i quits), then the best outcome under S′ i is no better than the worst outcome under Si. If S′ i(Ii) / ∈A0, then, since S′ i(Ii) ̸= Si(Ii), |A(Ii) \ A0| > 1. Then, by (1.c.iii) of Definition 18, ˜ t1 i will fall no further. So Si(Ii) guarantees that i will be in the allocation and receive transfer ˜ t1 i (Ii). But, by (1.d) of Definition 18, the best possible outcome under S′ i conditional on reaching Ii is no better, so the obvious dominance inequality holds.
Suppose θi + ˜ t1 i (Ii) ≤t0 i . Then, under Si, conditional on reaching Ii, agent i is not in the allocation and has transfer t0 i .53 However, under S′ i, either the outcome is the same, or agent i is in the allocation for some transfer ˆ ti ≤˜ t1 i (Ii) ≤t0 i −θi. Thus, the best possible outcome under S′ i is no better than the worst possible outcome under Si, and the obvious dominance inequality holds.
The argument proceeds symmetrically for Clause 2.
Notice that the above strategies result in some allocation and some payments, as a function of the type profile. We define these to be (fy, ft), such that G OSP-implements (fy, ft).
A.6 Theorem 5 Proof.
Lemma 6. For all B ⊆N, for all θN\B, ˜ ΘB(θN\B) is a join-semilattice with respect to the product order on ΘB.
Take any θ′′ B, θ′′′ B ∈˜ ΘB(θN\B). We want to show that θ′′ B ∨θ′′′ B ∈˜ ΘB(θN\B).
For all i ∈B, θ′′ B, θ′′′ B ∈closure({θB : ∀θ′ B\i ≥θB\i : i / ∈fy(θi, θ′ B\i, θN\B)}) (92) The set on the RHS is upward-closed with respect to the product order on θB\i.
Consider ˇ θB ≡θ′′ B ∨θ′′′ B. Its ith element has the property: ˇ θi = max{θ′′ i , θ′′′ i }. WLOG, suppose θ′′ i ≥θ′′′ i . Then, since ˇ θB\i ≥θ′′ B\i, ˇ θB ∈closure({θB : ∀θ′ B\i ≥θB\i : i / ∈fy(θi, θ′ B\i, θN\B)}) (93) Since the above argument holds for all i ∈B, θ′′ B ∨θ′′′ B ∈˜ ΘB(θN\B). This concludes the proof of Lemma 6.
First we prove the “if” direction. We do this by constructing G (and the corresponding strategy profiles). ft is specified implicitly.
Fix, for each i, the partition points {θk i }K k=1. Initialize k0 := (1, 1, . . . , 1), where k0 i denotes the ith element of this vector. Each agent i chooses whether to stay in the auction, at price θk0 i i . i 53By (1.c.ii) of Definition 18, either i will have quit in the past, or will have an opportunity to quit now, which he exercises.
46 quits iffi has a type weakly less than θk0 i i . Set B0 := N. (We use Bl to keep track of the agents still bidding in iteration l.) Set y0 := ∅. (We use yl to keep track of agents who must be in the allocation at iteration l.) At each stage, we define θQ N\Bl−1 ≡{θ kl−1 i i }i∈N\Bl−1. These are the recorded type (intervals) of the agents who have stopped bidding.
For l = 1, 2, . . .: 1. If Bl−1 = ∅, then terminate the algorithm at allocation y = yl−1.
2. If (θ kl−1 i i )i∈Bl−1 = sup{˜ ΘBl−1(θQ N\Bl−1)} (94) then (a) Choose agent i ∈Bl−1 such that, if θBl−1 > (θ kl−1 j j )j∈Bl−1, then i ∈fy(θBl−1, θQ N\Bl−1).
(b) Charge that agent the price θ kl−1 i i .
(c) Ask that agent to report ˆ k > kl−1 i such that θi ∈(θ ˆ k−1 i , θˆ k i ]. Set (kl j)j∈N such that: kl j := (ˆ k if j = i kl−1 j otherwise.
(95) (d) Set Bl := Bl−1 \ i (e) Set yl := yl−1 ∪i (f) Skip to stage l + 1.
3. Choose i ∈Bl−1 such that (kl j)j∈N satisfies kl j := ( kl−1 j + 1 if j = i kl−1 j otherwise.
(96) and (θ kl j j )j∈Bl−1 ∈˜ ΘBl−1(θQ N\Bl−1) (97) 4. Offer agent i the option to quit. Agent i quits iffhis type is less than or equal to θkl i i .
5. If agent i does not quit, set Bl := Bl−1.
6. If agent i quits, set Bl := Bl−1 \ i.
7. Set yl := yl−1 8. Go to stage l + 1.
The above algorithm defines an auction with monotonically ascending prices, where an agent has the option to quit (for a transfer of zero) whenever her price rises. i makes a payment equal to her going price at the first point where she ‘clinches the object’ - i.e. when she is guaranteed to 47 be in the final allocation. When i clinches the object, i is also asked to report her type - this may affect the payoffs of the agents that remain active, but does not affect her.54 By inspection, it is an obviously dominant strategy for i to stay in the auction if the price is less than her type, to quit if the price is greater than or equal to her type and to report her type truthfully at the point when she clinches the object.
It remains to show that the algorithm is well-defined for all type profiles, and that whenever it terminates, the resulting allocation agrees with fy. In particular, we must show that in Steps (2a) and (3), we can pick agent i satisfying the requirements of the algorithm.
Step (2a) is well-defined by assumption.
Lemma 7. (Given the assumptions of Theorem 5.) Under the above algorithm, for all l, (θkl i i )i∈Bl ∈ ˜ ΘBl(θQ N\Bl) We prove Lemma 7 by induction. It holds for l = 0 by the assumption that for all i, for all θ−i, i / ∈fy(θi, θ−i). Suppose it holds for l −1. We now prove that it holds for l (assuming, of course, that the algorithm does not terminate in Step 1 of iteration l).
Suppose (θ kl−1 i i )i∈Bl−1 = sup{˜ ΘBl−1(θQ N\Bl−1)}, so that Step 2 of the algorithm is triggered in iteration l. Since (θ kl−1 i i )i∈Bl−1 ∈˜ ΘBl−1(θQ N\Bl−1), we know that for all j ∈Bl−1: (θ kl−1 i i )i∈Bl−1 ∈closure({θBl−1 : ∀θ′ Bl−1\j ≥(θ kl−1 i i )i∈Bl−1\j : j / ∈fy(θ kl−1 j j , θ′ Bl−1\j, θQ N\Bl−1)}) (98) The set on the RHS of Equation 98 is upward closed with respect to the product order on θBl−1\j. Bl ⊂Bl−1 and (θkl i i )i∈Bl = (θ kl−1 i i )i∈Bl. Moreover, for the agent i who just clinched the object, θkl i i > θ kl−1 i i . Consequently, for all j ∈Bl (θkl i i )i∈Bl ∈closure({θBl : ∀θ′ Bl\j ≥(θkl i i )i∈Bl\j : j / ∈fy(θ kl j j , θ′ Bl\j, θQ N\Bl)}) (99) Since this holds for each set in the intersection that defines ˜ ΘBl(θQ N\Bl), this entails that (θkl i i )i∈Bl ∈˜ ΘBl(θQ N\Bl).
Suppose (θ kl−1 i i )i∈Bl−1 ̸= sup{˜ ΘBl−1(θQ N\Bl−1)}, so that we reach Step 3 of the algorithm in iteration l. Then, provided Step 3 is well-defined (i.e. we can pick i satisfying our requirements), (θ kl j j )j∈Bl−1 ∈˜ ΘBl−1(θQ N\Bl−1) (100) Bl ⊆Bl−1, and if i ∈N \ Bl, then θQ i = θkl i i . Thus, 54We could remove this feature by restricting attention to non-bossy allocation rules, where if changing i’s type changes the allocation, then it also changes whether i is satisfied. However, the canonical results for SP do not assume non-bossiness of fy, and we do not do so here. Alternatively, we could rule out such OSP mechanisms by instead requiring full implementation. However, the canonical monotonicity results for SP hold only for weak implementation.
48 (θ kl j j )j∈Bl ∈˜ ΘBl(θQ N\Bl) (101) So we need only show that Step 3 is well-defined for iteration l, given that Lemma 7 holds for l −1. This will simultaneously prove Lemma 7, and demonstrate that Step 3 is well-defined throughout.
We know that (θ kl−1 j j )j∈Bl−1 ∈˜ ΘBl−1(θQ N\Bl−1) (102) and (θ kl−1 j j )j∈Bl−1 ̸= sup{˜ ΘBl−1(θQ N\Bl−1)} (103) Let ˆ B ⊆Bl−1 be the set of all agents in Bl−1 such that θ kl−1 j j is less than the jth element of sup{˜ ΘBl−1(θQ N\Bl−1)}. Define θ′ ˆ B = 0.5(θ kl−1 j j )j∈Bl−1 + 0.5(θ kl−1 j +1 j )j∈Bl−1.
Now we define two disjoint open sets: ΘL Bl−1 = {θBl−1 : θ ˆ B < θ′ ˆ B} (104) ΘH Bl−1 = [closure(ΘL Bl−1)]C (105) The sets ˜ ΘBl−1(θQ N\Bl−1) ∩ΘL Bl−1 and ˜ ΘBl−1(θQ N\Bl−1) ∩ΘH Bl−1 are disjoint nonempty sets, and are open in the subspace topology. By connectedness, there exists some θ′′ Bl−1 ∈˜ ΘBl−1(θQ N\Bl−1) \ (ΘL Bl−1 ∪ΘH Bl−1). Fix some θ′′ Bl−1. For some dimension i ∈ˆ B, θ′′ i = (θ kl−1 i i + θ kl−1 i +1 i )/2.
Define θ′′′ Bl−1 = (θ kl−1 j j )j∈Bl−1 ∨θ′′ Bl−1. By Lemma 6, ˜ ΘBl−1(θQ N\Bl−1) is a join-semilattice. Thus, θ′′′ Bl−1 ∈˜ ΘBl−1(θQ N\Bl−1).
By construction, (θ kl−1 j j )j∈ˆ B ≤θ′′′ ˆ B < (θ kl−1 j +1 j )j∈ˆ B (106) θ′′′ Bl−1\ ˆ B = (θ kl−1 j j )j∈(Bl−1\ ˆ B) (107) Pick i ∈ˆ B such that θ′′ i = (θ kl−1 i i + θ kl−1 i +1 i )/2. Define: kl j := ( kl−1 j + 1 if j = i kl−1 j otherwise.
(108) Consider the sequence 1 t θ′′′ Bl−1 + (1 −1 t )(θ kl j j )j∈Bl−1 for t = 1, 2, 3, . . .. Since fy admits a finite partition, and this sequence does not move past any partition coordinates, every element of this sequence is in ˜ ΘBl−1(θQ N\Bl−1). Since ˜ ΘBl−1(θQ N\Bl−1) is closed, the limit of the sequence, (θ kl j j )j∈Bl−1, is also in ˜ ΘBl−1(θQ N\Bl−1). This proves Lemma 7.
Now we show that whenever the algorithm terminates, it agrees with fy.
49 By the assumption that for all i, for all θ−i, i / ∈fy(θi, θ−i), it follows that all the agents that quit at price θ1 i = θi are never satisfied (i.e. i / ∈fy(θ∗) for the true type profile θ∗).
By construction, after any iteration l, the bidders that remain active Bl have true types (θ∗ i )i∈Bl that strictly exceed the going prices (θkl i i )i∈Bl.
The bidders that are inactive have their types recorded (as accurately as we need given the finite partition) in the vector θQ N\Bl.
Suppose Step 1 is not activated and Step 2 is activated, in iteration l. Then, based on the information revealed up to iteration l −1, we know that the chosen bidder i is such that i ∈fy(θ∗).
Thus, for all l, yl ⊆fy(θ∗).
Suppose neither Step 1 nor Step 2 is activated in iteration l. By Lemma 7, (θ kl−1 j j )j∈Bl−1 ∈ ˜ ΘBl−1(θQ N\Bl−1). Consider the chosen bidder i whose price is incremented. The new price vector satisfies: kl j := ( kl−1 j + 1 if j = i kl−1 j otherwise.
(109) and (θ kl j j )j∈Bl−1 ∈˜ ΘBl−1(θQ N\Bl−1) (110) We now show that if bidder i quits, i.e. if his true type θi is in the interval (θ kl−1 i i , θkl i i ], then he is required not to be in the allocation based on the information revealed thus far. We proceed by contradiction: If there exists θ∗ Bl−1\i > (θ kl j j )j∈Bl−1\i such that i ∈fy(θkl i i , θ∗ Bl−1\i, θQ N\Bl−1), then (θ kl j j )j∈Bl−1 / ∈˜ ΘBl−1(θQ N\Bl−1). (Thus contradicting Equation 110.) Pick ϵ such that 0 < ϵ < min{θkl i i −θ kl−1 i i , minj∈Bl−1\i{θ∗ j −θ kl j j }}.
For all θ∗∗ Bl−1 such that ∥θ∗∗ Bl−1 −(θ kl j j )j∈Bl−1∥2 < ϵ, it follows (by the finite partition and by monotonicity) that i ∈ fy(θ∗∗ i , θ∗ Bl−1\i, θQ N\Bl−1), and moreover θ∗∗ Bl−1\i < θ∗ Bl−1\i. Thus, θ∗∗ Bl−1 / ∈{θBl−1 : ∀θ′ Bl−1\i ≥θBl−1\i : i / ∈fy(θi, θ′ Bl−1\i, θQ N\Bl−1}. Since this is true for all θ∗∗ Bl−1 in an ϵ-ball around (θ kl j j )j∈Bl−1, it fol-lows that (θ kl j j )j∈Bl−1 / ∈closure({θBl−1 : ∀θ′ Bl−1\i ≥θBl−1\i : i / ∈fy(θi, θ′ Bl−1\i, θQ N\Bl−1}). Thus (θ kl j j )j∈Bl−1 / ∈˜ ΘBl−1(θQ N\Bl−1). We have the desired contradiction.
Thus, whenever some bidder i’s going price rises in iteration l, the types who quit are those that, based on the information revealed so far, are required by the allocation rule not to be satisfied.
For all l, for the true type profile θ∗, (Bl)C ∩(yl)C ⊆fy(θ∗)C.
Gathering results: For all l, yl = (Bl)C ∩yl ⊆fy(θ∗) and (Bl)C ∩(yl)C ⊆fy(θ∗)C. Thus, whenever Bl = ∅, fy(θ∗) = yl. This completes the proof of the “if” direction.
Now the “only if” direction. G OSP-implements (fy, ft), so fy is SP-implementable. Thus, fy is monotone.
For all i, type θi is never satisfied, and always has a zero transfer. Thus, by Theorem 3, we can restrict our attention to monotone price mechanisms that satisfy the “Either” clause in Definition 18 - i.e. every agent faces an ascending price associated with being satisfied, and a fixed outside option (call this an ascending price mechanism, or APM). Suppose we have some G that OSP-implements (fy, ft). Moreover, suppose G is pruned, so that G is an APM.
Take any B ⊆N and θN\B. We now show that ˜ ΘB(θN\A) is connected. Let p : [0, 1] →ΘB be the price path under G faced by agents in B, when the type profile for the agents in B is 50 sup{˜ ΘB(θN\B)} and the type profile for the agents in N \ B is θN\B. Let z be the terminal history that results from that type profile, and let l be the number of elements of that sequence. Let h1, h2, . . . , hl be the subhistories of z.
Formally, p is defined as follows: Start f(0) = (θi)i∈A.
For each subhistory hm, let p( m l ) be equal to the prices faced by agents in B at hm.
For all points in r ∈( m−1 l , m l ), p(r) = (1 −β)p(m −1) + βp(m), for β = (r −m−1 l )/(1/l).
By inspection, p is a continuous function. Moreover, since at any point when an agent i quits under G, i / ∈fy(θ∗) based on the information revealed so far, for all r, p(r) ∈˜ ΘB(θN\B). Thus, p is a path from ˜ ΘB(θN\B) to sup{˜ ΘB(θN\A)}.
By Lemma 6, ˜ ΘB(θN\B) is a join semi-lattice.
We can generate a path p′ from any θ′ B ∈ ˜ ΘB(θN\B) to sup{˜ ΘB(θN\B)}, by defining p′(r) ≡θ′ B ∨p(r). Thus, ˜ ΘB(θN\B) is path-connected, which implies that it is connected.
We now show that there exists i ∈B such that, if θB > sup{˜ ΘB(θN\B)}, then i ∈fy(θB, θN\B).
If we cannot choose some θB > sup{˜ ΘB(θN\B)}, then this holds vacuously. Thus, fix some θ′ B > sup{˜ ΘB(θN\B)}.
Let z be the terminal history in G when the type profile is (θ′ B, θN\B).
Suppose there does not exist i ∈B such that such that, for all θB > sup{˜ ΘB(θN\B)}, i ∈ fy(θB, θN\B). By definition of sup{˜ ΘB(θN\B)}, there also does not exist i ∈B such that such that, for all θB > sup{˜ ΘB(θN\B)}, i / ∈fy(θB, θN\B).
Thus, the price path for agents in B along history z (defined as before) is not such that p(r) ≤ sup{˜ ΘB(θN\B)} for all r ∈[0, 1]. Thus, there must be a first point along z where the price path is not in {˜ ΘB(θN\B)}. Consider the agent i whose price was incremented at that point.
For all j ∈B, the relevant set in the intersection that defines {˜ ΘB(θN\B)} is upward-closed with respect to the product order on B \ j. Thus, when the price first leaves {˜ ΘB(θN\B)} at some subhistory ht, it must be that p(t l ) / ∈closure({θB : ∀θ′ B\i ≥θB\i : i / ∈fy(θi, θ′ B\i, θN\B)}) (111) The complement of the set on the RHS of Equation 111 is open. Thus, for some ϵ > 0, an open ϵ-ball around p( t l) is a subset of the complement of the RHS.
Consequently, we can choose some θ′′ i strictly greater than i’s old price, and strictly less than i’s new going price, and some θ′′ B\i strictly greater (in the product order) than the going prices for B \ i, such that i ∈fy(θ′′ i , θ′′ B\i, θN\B). Since G is an APM, the actions of types (θ′′ i , θ′′ B\i) and the actions of types θ′ B are indistinguishable prior to that point. Thus, G does not result in the prescribed outcome for type profile (θ′′ i , θ′′ B\i, θN\B), a contradiction.
This completes the proof of the “only-if” direction.
B Alternative Empirical Specifications Here we report alternative empirical specifications for the experiment. Table 5 and 6 are identical to Table 2 except that they compute p-values and standard errors using alternative methods.
A natural measure of errors would be to take the sum, for k = 1, 2, 3, 4, of the absolute difference between the kth highest bid and the kth highest value. However we do not observe the highest bid under AC, and we often do not observe the highest bid under AC+X. We could instead take the 51 Table 5: mean(abs(2nd bid - 2nd value), p-values calculated using Wilcoxon rank-sum test Format Rounds SP OSP p-value Auction 1-5 8.04 3.19 < .001 6-10 4.99 1.77 .005 +X Auction 1-5 3.99 1.83 .004 6-10 3.69 1.29 < .001 This is the same as Table 2, except that the p-values are calculated using the Wilcoxon rank-sum test.
Table 6: mean(abs(2nd bid - 2nd value), p-values and standard errors calculated using clustered regression (clustered by groups) Format Rounds SP OSP p-value Auction 1-5 8.04 3.19 .006 (1.26) (1.06) 6-10 4.99 1.77 .014 (1.20) (0.34) +X Auction 1-5 3.99 1.83 .005 (0.61) (0.41) 6-10 3.69 1.29 .015 (0.88) (0.33) This is the same as Table 2, except that the p-values and standard errors are calculated by running a single regression (with appropriate indicator variables) and clustering by groups.
52 sum for k = 2, 3, 4 of the absolute difference between the kth highest bid and the kth highest value, averaged as before in five-round blocks. Table 7 reports the results.
Table 7: mean(sum(abs(kth bid - kth value))), for k = 2, 3, 4 Format Rounds SP OSP p-value Auction 1-5 32.63 9.89 < .001 (4.64) (1.89) 6-10 16.28 5.53 .001 (2.73) (0.91) +X Auction 1-5 17.04 6.18 .011 (3.70) (1.06) 6-10 14.21 4.74 .022 (3.70) (0.75) For each auction, we sum the absolute differences between the kth bid and the kth value, for k = 2, 3, 4. We then take the mean of this over each 5-round block. We then compute standard errors counting each group’s 5-round mean as a single observation. (18 observations per cell.) p-values are computed using a two-sample t-test, allowing for unequal variances.
Another measure of errors would be to take the sum of the absolute difference between each bidder’s bid and that bidder’s value, dropping all highest bidders for symmetry. Table 8 reports the results.
Table 9 reports the results of Table 3, except that the p-values are calculated using the Wilcoxon rank-sum test.
29.0% of rank-order lists are incorrect under SP-RSD. 2.6% of choices are incorrect under OSP-RSD. However, this is not a fair comparison; rank-order lists mechanically allow us to spot more errors than single choices. To compare like with like, we compute the proportion of incorrect choices we would have observed, if subjects played OSP-RSD as though they were implementing the submitted rank-order lists for SP-RSD. This is a cautious measure; it counts errors under SP-RSD only if they would have altered the outcome under OSP-RSD. Table 10 reports the results.
Finally, in Table 11, we compute the mean difference between the second-highest bid and the second-highest value, by auction format and by five-round blocks. This summarizes the average direction of deviations (positive or negative) from equilibrium play. I had no prior hypothesis about these outcome variables, but report this analysis for completeness.
C Quantal Response Equilibrium Quantal response equilibrium (QRE) is defined for normal form games (McKelvey and Palfrey, 1995). Agent quantal response equilibrium (AQRE) adapts QRE to extensive form games (McK-elvey and Palfrey, 1998). In AQRE, each agent plays a (perturbed) best response at each informa-tion set, given correct beliefs about other agents’ play and his own mistakes at future information sets. In AC and AC+X, a rational agent is indifferent between quitting and continuing when the price is equal to his value, since if he continues presently, then he will quit at the next bid increment.
Consequently, an AQRE agent strictly prefers to quit when the price is equal to his value, since he correctly forsees that if he continues, he will sometimes make a mistake at the next bid increment, and the increment after that, resulting in a negative expected value. By continuity, an AQRE agent strictly prefers to quit when the price is in some interval below his value. Consequently, AQRE 53 Table 8: mean(sum(abs(i’s bid - i’s value))), dropping highest bidders Format Rounds SP OSP p-value Auction 1-5 35.13 10.18 < .001 (5.20) (1.88) 6-10 15.46 4.89 .002 (2.85) (0.72) +X Auction 1-5 17.88 5.58 .009 (4.10) (1.01) 6-10 14.20 4.64 .022 (3.72) (0.83) For each auction, we sum the absolute differences between each bidder’s bid and their value, dropping the highest bidder. We then take the mean of this over each 5-round block. We then compute standard errors counting each group’s 5-round mean as a single observation. (18 observations per cell.) p-values are computed using a two-sample t-test, allowing for unequal variances.
Table 9: Proportion of serial dictatorships not ending in dominant strategy outcome, p-values calculated using Wilcoxon rank-sum test SP OSP p-value Rounds 1-5 43.3% 7.8% .0001 Rounds 6-10 28.9% 6.7% .0010 This is the same as Table 3, except that the p-values are calculated using the Wilcoxon rank-sum test.
Table 10: Proportion of incorrect choices under serial dictatorship: SP (imputed) vs OSP (actual) SP OSP p-value Rounds 1-5 17.8% 2.6% < .001 (3.5%) (1.1%) Rounds 6-10 10.7% 2.6% .002 (2.0%) (1.3%) p-value .078 1.000 For each group in SP-RSD, for each period, we simulate the three choices that we would have observed under OSP-RSD. For each group, for each 5-round block, we record the proportion of choices that are incorrect. We then compute standard errors counting each group-block pair as a single observation. (18 observations per cell.) When comparing SP to OSP, we compute p-values using a two-sample t-test. When comparing early to late rounds of the same game, we compute p-values using a paired t-test. In the sample for OSP-RSD, there are 7 incorrect choices in the first five rounds and 7 incorrect choices in the last five rounds.
54 predicts systematic under-bidding in AC and AC+X.
As Figure 2 and Figure 3 illustrate, there is no evidence of systematic under-bidding in AC and AC+X. Table 11 shows this rigorously.
D Experiment instructions 55 WELCOME This is a study about decision-making. Money earned will be paid to you in cash at the end of the experiment. This study is about 90 minutes long.
We will pay you $5 for showing up, and $15 for completing the experiment. Additionally, you will be paid in cash your earnings from the experiment. If you make choices in this experiment that lose money, we will deduct this from your total payment. However, your total payment (including your show-up payment and completion payment) will always be at least $20.
You have been randomly assigned into groups of 4. This experiment involves 3 games played for real money. You will play each game 10 times with the other people in your group. We will give you instructions about each game just before you begin to play it. Your choices in one game will not affect what happens in other games.
There is no deception in this experiment. Every game will be exactly as specified in the instructions. Anything else would violate the IRB protocol under which we run this study. (IRB Protocol 34876) Please do not use electronic devices or talk with other volunteers during this study. If we do find you using electronic devices or talking with other volunteers, the rules of the study require us to deduct $20 from your earnings.
If you have questions at any point, please raise your hand and we will answer your questions privately. # of # 1 6 BB GAME 1 In this game, you will bid in an auction for a money prize. The prize may have a different dollar value for each person in your group. You will play this game for 10 rounds. All dollar amounts in this game are in 25 cent increments.
At the start of each round, we display your value for this round’s prize. If you win the prize, you will earn the value of the prize, minus any payments from the auction. Your value for the prize will be calculated as follows: 1. For each group we will draw a common value, which will be between $10.00 and $100.00. Every number between $10.00 and $100.00 is equally likely to be drawn.
2. For each person, we will also draw a private adjustment, which will be between $0.00 and $20.00. Every number between $0.00 and $20.00 is equally likely to be drawn.
In each round, your value for the prize is equal to the common value plus your private adjustment. At the start of each round, you will learn your total value for the prize, but not the common value or the private adjustment.
This means that each person in your group may have a different value for the prize. However, when you have a high value, it is more likely that other people in your group have a high value.
The auction proceeds as follows: First, you will learn your value for the prize. Then you can choose a bid in the auction. Each person in your group will submit their bids privately and at the same time. You do this by typing your bid into a text box and clicking ‘confirm bid’. You will have 90 seconds to make your decision, and can revise your bid as many times as you like. At the end of 90 seconds, your final bid will be the one that counts.
of # 2 6 BB All bids must be between $0.00 and $150.00, and in 25 cent increments. The highest bidder will win the prize, and make a payment equal to the second-highest bid. This means that we will add to her earnings her value for the prize, and subtract from her earnings the second-highest bid. All other bidders’ earnings will not change.
At the end of each auction, we will show you the bids, ranked from highest to lowest, and the winning bidder’s profits. If there is a tie for the highest bidder, no bidder will win the object. # of # 3 6 BB GAME 2 In this game, you will bid in an auction for a money prize. You will play this game for 10 rounds.
Your value for the prize will be generated as before.
However, each round, we will also draw a new number, X, for each group. The rules of the auction are different, as follows: All bidders will submit their bids privately and at once. However, the highest bidder will win the prize if and only if their bid exceeds the second-highest bid by more than X.
If the highest bidder wins the prize, she will make a payment equal to the second-highest bid plus X. This means that we will add to her earnings her value for the prize, and subtract from her earnings the second-highest bid plus X. All other bidders’ earnings will not change.
If the highest bid does not exceed the second-highest bid by more than X, then no bidder will win the prize. In that case, no bidder’s earnings will change.
X will be between $0.00 and $3.00, with every 25 cent increment equally likely to be drawn. You will be told your value for the prize at the start of each round, but will not be told X. At the end of each round, we will tell you the value of X. # of # 4 6 BB GAME 3 You will play this game for 10 rounds. In each round of this game, there are four prizes, labeled A, B, C, and D. Prizes will be worth between $0.00 and $1.25. For each prize, its value will be the same for all the players in your group. At the start of each round, you will learn the value of each prize. You will also learn your priority score, which is a random number between 1 and 10. Every whole number between 1 and 10 is equally likely to be chosen.
The game proceeds as follows: We will ask you to list the prizes, in any order of your choice. All players will submit their lists privately and at the same time.
After all the lists have been submitted, we will assign prizes using the following rule: 1.
The player with the highest priority score will be assigned the top prize on his list.
2.
The player with the second-highest priority score will be assigned the top prize on his list, among the prizes that remain.
3.
The player with the third-highest priority score will be assigned the top prize on his list, among the prizes that remain.
4.
The player with the lowest priority score will be assigned whatever prize remains.
If two players have the same priority score, we will break the tie randomly.
You will have 90 seconds to form your list. You do this by typing a number, from 1 to 4, next to each prize, and then clicking the button that says “Confirm Choices”. Each prize must be assigned a different number, from 1 (top) to 4 (bottom). Your choices will not count unless you click the button that says “Confirm Choices”.
of # 5 6 BB If you do not produce a list by the end of 90 seconds, we will assign prizes as though you reported the list in order A-B-C-D.
At the end of each round, we will add to your earnings the value of the prize you were assigned.
of # 6 6 BB WELCOME This is a study about decision-making. Money earned will be paid to you in cash at the end of the experiment. This study is about 90 minutes long.
We will pay you $5 for showing up, and $15 for completing the experiment. Additionally, you will be paid in cash your earnings from the experiment. If you make choices in this experiment that lose money, we will deduct this from your total payment. However, your total payment (including your show-up payment and completion payment) will always be at least $20.
You have been randomly assigned into groups of 4. This experiment involves 3 games played for real money. You will play each game 10 times with the other people in your group. We will give you instructions about each game just before you begin to play it. Your choices in one game will not affect what happens in other games.
There is no deception in this experiment. Every game will be exactly as specified in the instructions. Anything else would violate the IRB protocol under which we run this study. (IRB Protocol 34876) Please do not use electronic devices or talk with other volunteers during this study. If we do find you using electronic devices or talking with other volunteers, the rules of the study require us to deduct $20 from your earnings.
If you have questions at any point, please raise your hand and we will answer your questions privately. # of # 1 5 AA GAME 1 In this game, you will bid in an auction for a money prize. The prize may have a different dollar value for each person in your group. You will play this game for 10 rounds. All dollar amounts in this game are in 25 cent increments.
At the start of each round, we display your value for this round’s prize. If you win the prize, you will earn the value of the prize, minus any payments from the auction. Your value for the prize will be calculated as follows: 1. For each group we will draw a common value, which will be between $10.00 and $100.00. Every number between $10.00 and $100.00 is equally likely to be drawn.
2. For each person, we will also draw a private adjustment, which will be between $0.00 and $20.00. Every number between $0.00 and $20.00 is equally likely to be drawn.
In each round, your value for the prize is equal to the common value plus your private adjustment. At the start of each round, you will learn your total value for the prize, but not the common value or the private adjustment.
This means that each person in your group may have a different value for the prize. However, when you have a high value, it is more likely that other people in your group have a high value.
The auction proceeds as follows: First, you will learn your value for the prize. Then, the auction will start. We will display a price to everyone in your group, that starts low and counts upwards in 25 cent increments, up to a maximum of $150.00. At any point, you can choose to leave the auction, by clicking the button that says “Stop Bidding”.
of # 2 5 AA When there is only one bidder left in the auction, that bidder will win the prize at the current price. This means that we will add to her earnings her value for the prize, and subtract from her earnings the current price. All other bidders’ earnings will not change.
At the end of each auction, we will show you the prices where bidders stopped, and the winning bidder’s profits. If there is a tie for the highest bidder, no bidder will win the object. # of # 3 5 AA GAME 2 In this game, you will bid in an auction for a money prize. You will play this game for 10 rounds.
Your value for the prize will be generated as before.
However, each round, we will also draw a new number, X, for each group.
The rules of the auction are different, as follows: The price will count up from a low value, and you can choose to leave the auction at any point, by clicking the button that says “Stop Bidding”. When there is only one bidder left in the auction, the price will continue to rise for another X dollars, and then freeze.
If the last bidder stays in the auction until the price freezes, then she will win the prize at the final price. This means that we will add to her earnings her value for the prize, and subtract from her earnings the final price. All other bidders’ earnings will not change.
If the last bidder stops bidding before the price freezes, then no bidder will win the prize. In that case, no bidder’s earnings will change.
X will be between $0.00 and $3.00, with every 25 cent increment equally likely to be drawn. You will be told your value for the prize at the start of each round, but will not be told X. At the end of each round, we will tell you the value of X. # of # 4 5 AA GAME 3 You will play this game for 10 rounds. In each round of this game, there are four prizes, labeled A, B, C, and D. Prizes will be worth between $0.00 and $1.25. For each prize, its value will be the same for all the players in your group.
At the start of each round, you will learn the value of each prize. You will also learn your priority score, which is a random number between 1 and 10. Every whole number between 1 and 10 is equally likely to be chosen.
The game proceeds as follows: 1.
The player with the highest priority score will pick one prize. 2.
The player with the second-highest priority score will pick one of the prizes that remains. 3.
The player with the third-highest priority score will pick one of the prizes that remains. 4.
The player with the lowest priority score will be assigned whatever prize remains.
If two players have the same priority score, we will break the tie randomly.
When it is your turn to pick, you will have 30 seconds to make your choice. You do this by selecting a prize and then clicking the button that says “Confirm Choice”. Your choices will not count unless you click the button that says “Confirm Choice”.
If you do not make a choice by the end of 30 seconds, we will assign prizes as though you picked whichever prize is earliest in the alphabet.
At the end of each round, we will add to your earnings the value of the prize you picked.
of # 5 5 AA Table 11: mean(2nd bid - 2nd value) Format Rounds SP OSP Auction 1-5 -1.59 0.58 (2.07) (0.77) 6-10 4.22∗∗∗ 0.78∗ (1.28) (0.40) +X Auction 1-5 1.62∗∗ 0.15 (0.70) (0.47) 6-10 2.97∗∗∗ 0.02 (0.97) (0.31) For each group, we take the mean difference between the second-highest bid and the second-highest value over each 5-round block. We then compute standard errors counting each group’s 5-round mean as a single observation. (18 observations per cell, standard errors in parentheses.) We then run a two-sided t-test for the null that the value in each cell has zero mean. denotes p < .10, denotes p < .05, and denotes p < .01.
67 |
13788 | https://www.quora.com/What-is-a-wing-strategy-for-the-game-nim-Can-you-prove-that-it-is-a-winning-strategy | We mathematicians have an entire area of mathematics called Combinatorial Game Theory (Combinatorial game theory - Wikipedia) that deals with games like Nim. And Nim is used as the iconic example, so much that the Grundy values in the Sprague–Grundy Theorem (Sprague–Grundy theorem - Wikipedia) are also known as nim-values or nimbers.
In general, such games can be won by using the Minimum Excluded Value strategy, called Mex for short (Mex (mathematics) - Wikipedia). But Nim has hidden symmetries that make it easier, so let me explain Nim alone.
Nim can be played with chalk tally marks on a blackboard, matchsticks in rows, stones in piles, poker chips in stacks, anything that lets us have separate stacks from which we can remove or erase items. The rule is that a player can remove any number of objects from a single stack during their turn, they must remove at least one object, and the person who removed the last object from the last stack wins. For example, a typical Nim game could start with three piles of 1 stone, 2 stones, and 3 stones. The first player removes both stones from the 2-stone pile. The second player removes 2 stones from the 3-stone pile. This leaves two piles with one stone each. The first player is stuck removing one pile and the second player removes the other to win.
One obvious Nim strategy is that if the game starts with two piles of the same size, then the second player can always win by mirroring the first player's move on the other pile. Thus, the second player always restores the game to having two equal piles. Eventually, the first player takes the last stone in one pile and the second player wins by taking the last stone in the other pile. The state of two equal piles is known as a losing position because the other player can win. And the state of two unequal piles is known as a winning position because the player can equalize the piles to leave the other player with two equal piles.
Two pairs of equal piles is also a losing position, because the second player just has to keep the two piles in each pair equal to each other.
Three piles of 1 stone, 2 stones, and 3 stones is a losing position, because whatever the first player does, the second player can leave behind two identical piles. I could demonstrate this by examining the six possible first moves: 1 stone from the 1-stone pile, 1 or 2 stones from the 2-stone pile, or 1, 2, or 3 stones from the 3-stone pile. To summarize, the three possible moves of removing an entire pile leave behind two unequal piles, a winning position for the second player and the other three possible moves of removing so few stones from a big pile to leave three piles behind leaves behind two equal piles and an unequal pile. The second player removes the unequal pile to leave behind two equal piles.
The basic principle is that if a position has at least one way of leaving behind a losing position for the opponent, such as two equal piles, then that position is winning. If all moves from a position leave behind a winning position for the opponent, then that position is winning.
Nim has a special trick that the winning positions can be identified by a mathematical formula that comes from writing the size of every pile in binary notation. The trick is a generalization on matching up equal-size piles.
Let's suppose we have four piles, a 3-stone pile, a 4-stone pile, a 5-stone pile, and a 6-stone pile. Write 3, 4, 5, and 6 in binary as 11, 100, 101, and 110. Then count their digits in each place that are 1's. 11 and 101 both have a 1 in the 1's place, so the 1's place is even. 11 and 110 both have a 1 in the 10's place, so the 10's place is even. 100, 101, and 110 all have a 1 in the 100's place, so the 100's place is odd. If any of the places have an odd number of 1's then the position is winning. If all of the places have an even number of 1's, then the position is winning.
The first player has a winning position, because he has an odd number of 1's in the 100's place. That player's move needs to remove a 1 in the 100's place without adding a 1 in any other place. The easiest way would be to remove all 4 stones (100 stones in binary) from the 4-stone pile, but removing 4 stones from the 5-stone or 6-stone pile would work just as good. Let me remove 4 stones from the 5-stone pile. That leaves a 1-stone pile, 3-stone pile, 4-stone pile, and 6-stone pile.
Suppose the second player then removes 1 stone from the 6-stone pile. That leaves a 1-stone pile, 3-stone pile, 4-stone pile, and 5-stone pile. In binary that is 1, 11, 100, and 101. The 100's place has even 1, but the 10's place and 1's place have odd ones, so we need to correct that. One way would be to remove all the stones from the 3-stone (11 in binary) pile. The first player cannot create another 1 in the 10's place without losing a 1 in the 100's place, so that is the only winning move. This leaves a 1-stone pile, 4-stone pile, and 5-stone pile for the second player.
Suppose the second player removed 1 stone from the 4-stone pile. This leaves a 1-stone pile, 3-stone pile, and 5-stone pile for the first player. In binary those numbers are 1, 11, and 101. All the places have an odd number of 1's, so all places need to change. Removing stones from the 5-stone pile is the only way to change the 100's place, so remove 3 stones from it to leave 2 stones. That gives the known losing position of a 1-stone pile, 2-stone pile, and 3-stone pile.
Why this method works is it is secretly pairing subpiles inside the piles. An even number of 1's in a place means a secret pairing. But an odd number of 1's means unequal piles. I don't know why the hidden subpiles have to be in powers-of-2 sizes, which are revealed in binary notation, but that is how it works.
The proof this strategy works is that having two equal piles always have an even number of 1’s in each binary place, and balancing those 1′s means keeping the piles equal. Thus, the binary strategy always wins in the endgame. For earlier in the game more piles, forcing the opponent to always start with an even number of 1’s in each binary place eventually forces the game into the two-equal-piles mode as the number of stones shrinks. And a player can always move from odd 1’s in some place to even 1’s in all places, so the winning player won’t get stuck. |
13789 | http://www.euonym.us/~academics/spring06/chm106/docs/lab6.pdf | CHM 106 Potentiometric Titration of Phosphoric Acid BACKGROUND Potentiometric titrations are a useful method of determining unknown concentrations in many different types of chemical systems. They may be employed in redox chemistry by monitoring the reduction potential of the titration flask as titrant is added, or they may be used in acid-base chemistry where the pH of the titration flask is monitored as acid or base is added. Experimentally, a pH meter relates a voltage from the probe to pH, so these techniques are actually quite similar. A detailed examination of potentiometric acid-base titrations is given in section 15.4 (pages 729-744) of the textbook. In this lab, we will conduct a potentiometric titration of a solution of phosphoric acid (H3PO4) of unknown concentration. We will determine the concentration of the phosphoric acid solution and the values for the first two acid dissociation constants. We will also examine indicator selection in titrations. MATERIALS AND REAGENTS 250 mL beaker, 10 mL pipet, buret, pH meter, magnetic stirplate and stirbar 0.100 M NaOH, H3PO4 unknown, pH 4 and 7 reference solutions, indicators PROCEDURE This lab will be conducted in groups but the lab report must be completed and submitted individually. Pipet 10.00 mL of the H3PO4 solution into a clean, dry 250 mL beaker. Add 40 mL of distilled water to the beaker. Place the 250 mL beaker on a magnetic stirrer and add a stirring bar. Fill a clean buret with 0.100 M NaOH solution so that it is exactly at the 0.00 mL mark and position over the beaker. Calibrate a pH meter using reference buffers of pH 4.0 and pH 7.0. Using an apparatus clamp, position the pH electrode in the 250 mL beaker so that the probe is submerged in the solution but is not in any danger of being struck by the stirbar. Begin stirring the solution at a medium speed. Add a couple drops of indicator solution. Each group will use a different indicator. Record the pH of the solution. Add about 0.5 mL of NaOH from the buret. Record accurately the volume of base added and the pH of the solution. Repeat, adding base in 0.5 mL increments and recording the pH until all 50 mL of base have been added. While adding base, observe the color of the indicator in your titration flask and record the endpoint of the titration, according to your indicator, along with the appropriate data points. Be sure to obtain the observations on the indicator color change from other groups before you leave the laboratory. DATA ANALYSIS Using spreadsheet software, plot the pH of the titration flask versus the volume of NaOH added. You should obtain a plot similar to the one shown below, but your values will be different because the concentrations and amounts used in the lab are different from this example. Sample Plot Titration of 1M H3PO4 with 1M KOH 0 2 4 6 8 10 12 14 0.00 0.02 0.04 0.06 0.08 0.10 0.12 0.14 0.16 0.18 0.20 Volume KOH (L) pH pH = pKa1 First equivalence point pH = pKa2 Second equivalence point Since phosphoric acid is a polyprotic acid, we will observe one equivalence point for each successive ionization. However, since the Ka for this third ionization (1.259x10-12) is close in magnitude to the ion product constant of water (1.00x10-14) the third equivalence point will be obscured on the titration curve. We will recall that the steep vertical regions correspond to the equivalence points and at half the volume of an equivalence point (or the volume halfway in between equivalence points) is the point at which the pH of the solution is equal to the pKa for that dissociation. However, it is difficult to determine exactly where on the vertical region the equivalence point lies. Observe that near the equivalence point, the slope of the titration curve is very high. Prior to the equivalence point, the slope of the titration curve is increasing, and after the equivalence point, the slope begins to decrease. Mathematically, the point at which the slope of a curve stops increasing and begins to decrease is called an inflection point. We can take advantage of mathematics to make it easier to determine precisely where the equivalence point lies. In your spreadsheet, calculate the slope between each pair of data points. Mathematically, V ) pH ( V V pH pH 1 2 1 2 d d slope = − − = Where pH1 and pH2 are the pH values of successive data points and V1 and V2 are the corresponding volumes of NaOH. This slope of the curve that we have calculated is called the first derivative and is designated as V ) pH ( d d because it is the rate of change of the pH with respect to volume. Prepare a plot of the first derivative of pH (slope of pH vs V) versus volume of NaOH added. You should obtain a plot similar to the one shown below: Sample Plot Titration of 1M H3PO4 with 1M KOH 0 100 200 300 400 500 600 700 800 900 0.00 0.02 0.04 0.06 0.08 0.10 0.12 0.14 0.16 0.18 0.20 Volume KOH (L) d(pH) / dV First equivalence point Second equivalence point Calculus tells us that when a function passes through an inflection point, its derivative will pass through a local maximum or minimum. Notice that it is considerably easier to find the locations of the equivalence points in the plot of the first derivative. However, due to experimental errors in the data we have, the tops of the peaks in this plot may be broadened, and still somewhat difficult to locate the precise maximum. We can once again turn to mathematics for help. Observe that the slope of the derivative plot is positive before reaching the maximum because the value of the derivative is increasing. At the maximum, the slope of the derivative is zero, and past the maximum the slope of the derivative plot is negative because the value of the derivative is decreasing. We can thus examine the derivative (slope) of the derivative plot and when it equals zero, we will know precisely the location of the equivalence point and the corresponding volume of NaOH at the equivalence point. In your spreadsheet, calculate the slope of the derivative between each set of data points. Mathematically, 2 2 1 2 1 2 V ) pH ( V V ' pH ' pH d d slope = − − = where pH’1 and pH’2 are successive values of the derivative (not the pH itself) and V1 and V2 are the corresponding volumes of NaOH. You should obtain a plot similar to the following: Sample Plot Titration of 1M H3PO4 with 1M KOH -350000 -250000 -150000 -50000 50000 150000 250000 350000 0.00 0.02 0.04 0.06 0.08 0.10 0.12 0.14 0.16 0.18 0.20 Volume KOH (L) d2(pH) / dV2 First equivalence point Second equivalence point Now it is unambiguous where the equivalence point lies because the is a very sharp transition of the second derivative plot from positive to negative and an obvious value where it equals zero. This value is the volume of NaOH at the equivalence point. From your second derivative plot, determine the volume of sodium hydroxide at the first and second equivalence point. Once we know the volume of sodium hydroxide at each equivalence point, it is a simple matter to find the pKa values for each dissociation. Refer to the first plot you constructed. At a volume equal to half the equivalence point (or halfway between equivalence points), pH = pKa. From your first plot, determine the values of pKa1 and pKa2 for phosphoric acid. Given that the literature values are pKa1 = 2.148 and pKa2 = 7.198, calculate a percentage error for each measurement. We can also use this data to determine the concentration of H3PO4. At the first equivalence point, the reaction ) ( PO H ) ( O H ) ( OH ) ( PO H 4 2 2 4 3 aq l aq aq − − + → + is complete. Since we know the concentration (0.100 M) of NaOH, the volume of NaOH at the equivalence point, and the volume of H3PO4 used, we can calculate the concentration of phosphoric acid. Similarly, at the second equivalence point, the reaction ) ( HPO ) ( O H 2 ) ( OH 2 ) ( PO H 2 4 2 4 3 aq aq aq aq − − + → + is complete so we can do similar calculations. These calculations should result in the same value, but they will probably be slightly different due to experimental error. Determine the value of the concentration of phosphoric acid using data from each equivalence point. Average their values and report this average as the value of the concentration of phosphoric acid. Mark the end point of each titration, according to the indicators used, on your plot of pH versus volume of NaOH. In your discussion, comment on whether each indicator was appropriate for the titration (Was the end point close to the equivalence point?) and comment on the implications of the end point of each indicator on doing the calculations for the titration. Calculate the concentration of phosphoric acid according to each indicator’s end point, and calculate a percentage error between the end point and the appropriate equivalence point. QUESTIONS 1. Would bromothymol blue (pKa = 7.1) be an appropriate indicator for the titration of phosphoric acid samples? Explain. 2. How does buffering in a titration of a weak acid or weak base affect the shape of the titration curve when compared to the titration curve of strong acid with strong base? Explain. |
13790 | https://www.merckmillipore.com/SR/en/product/mm/818755 | Acetic acid 99-100 for synthesis 64-19-7
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Acetic acid
99-100%, for synthesis
Synonym(s):
Acetic acid, Ethanoic acid, Glacial acetic acid
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About This Item
Linear Formula:
CH 3 CO 2 H
CAS Number:
64-19-7
Molecular Weight:
60.05
Beilstein:
506007
MDL number:
MFCD00036152
UNSPSC Code:
12352106
EC Index Number:
200-580-7
NACRES:
NA.22
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Properties
Product Name
Acetic acid 99-100%, for synthesis
grade
synthesis grade
Quality Level
200
vapor density
2.07 (vs air)
vapor pressure
15.4 hPa ( 20 °C)
Assay
≥99% (acidimetric)
form
liquid
autoignition temp.
485°C
800°F
potency
3310 mg/kg LD 50, oral (Rat)
expl. lim.
16%, 92°F
4%, 59°F
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| ###### This Item | 1.00062 | 5.43808 | 1.59166 |
--- --- |
| Sigma-Aldrich 8.18755 Acetic acid 99-100% Quick View | Supelco 1.00062 Acetic acid 96% Quick View | Supelco 5.43808 Acetic acid 100% Quick View | Supelco 1.59166 Acetic acid 30% Quick View |
| assay ≥99% (acidimetric) | assay ≥96.0% (alkalimetric) | assay ≥99.8% (alkalimetric) | assay - |
| Quality Level 200 | Quality Level 300 | Quality Level 100 | Quality Level 300 |
| form liquid | form liquid | form liquid | form liquid |
| storage temp. 15-25°C | storage temp. 15-25°C | storage temp. 15-25°C | storage temp. no temp limit |
| solubility soluble 602.9 g/L | solubility - | solubility - | solubility - |
| pH 2.5 (20°C, 50 g/L in H 2 O) | pH 2.5 (20°C, 10 g/L in H 2 O) | pH 2.5 (20°C, 50 g/L in H 2 O) | pH - |
Description
Analysis Note
Assay (acidimetric): ≥ 99.0 %
Identity: passes test
Density (d 20/4): 1.048 - 1.052
Safety Information
Pictograms
GHS02,GHS05
Signal Word
Danger
Hazard Statements
H226 - H314
Precautionary Statements
P210 - P233 - P240 - P280 - P303 + P361 + P353 - P305 + P351 + P338
Hazard Classifications
Eye Dam. 1 - Flam. Liq. 3 - Skin Corr. 1A
Storage Class Code
3 - Flammable liquids
WGK
WGK 1
Flash Point(F)
102.2 °F
Flash Point(C)
39 °C
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13791 | https://en.wiktionary.org/wiki/%E9%80%99%E5%80%8B | Jump to content
這個
한국어
Magyar
Malagasy
日本語
中文
From Wiktionary, the free dictionary
See also: 这个
Chinese
[edit]
| | | this | (a measure word); individual |
--- --- |
| trad. (這個) | | 這 | 個 |
| simp. (这个) | | 这 | 个 |
| alternative forms | | 這个 Hokkien 即個/即个 Hokkien 此個/此个 Hokkien | |
Pronunciation
[edit]
Mandarin
: (Standard)
: (Pinyin): zhège, zhèige
: (Zhuyin): ㄓㄜˋ ˙ㄍㄜ, ㄓㄟˋ ˙ㄍㄜ
: Duration: 2 seconds.0:02
: (Dungan, Cyrillic and Wiktionary): җыгә (žɨgə, III-I)
Cantonese (Jyutping): ze2 go3 / ze5 go3
Hakka
: (Sixian, PFS): liá-ke / yá-ke
: (Meixian, Guangdong): zê3 gê4 / ê3 gê4 / gê3 gê4
Southern Min (Hokkien, POJ): chit-ê / chit-gê / chit-gé / chiak-gê / chit-e / chih-ge / chih-gê / chek-e / chek-ê / chek-ge / chek-gê / chit-le / chit-lê / chí-ê / chí-gé / chit-é / chih-lê / chí-lê / che / chémore ▼
Mandarin
(Standard Chinese)+
Hanyu Pinyin: zhège
Zhuyin: ㄓㄜˋ ˙ㄍㄜ
Tongyong Pinyin: jhège̊
Wade–Giles: chê4-ko5
Yale: jè-ge
Gwoyeu Romatzyh: jeh.ge
Palladius: чжэгэ (čžɛgɛ)
Sinological IPA (key): /ʈ͡ʂɤ⁵¹ g̊ə¹/
(Standard Chinese)+
Hanyu Pinyin: zhèige
Zhuyin: ㄓㄟˋ ˙ㄍㄜ
Tongyong Pinyin: jhèige̊
Wade–Giles: chei4-ko5
Yale: jèi-ge
Gwoyeu Romatzyh: jey.ge
Palladius: чжэйгэ (čžɛjgɛ)
Sinological IPA (key): /ʈ͡ʂeɪ̯⁵¹ g̊ə¹/
Note: zhèige is a coalescence of 這一個 / 这一个.
(Dungan)
Cyrillic and Wiktionary: җыгә (žɨgə, III-I)
Sinological IPA (key): /ʈ͡ʂʐ̩⁴⁴ kə²⁴/
: (Note: Dungan pronunciation is currently experimental and may be inaccurate.)
Cantonese
+ (Standard Cantonese, Guangzhou–Hong Kong)+
- Jyutping: ze2 go3 / ze5 go3
- Yale: jé go / jéh go
- Cantonese Pinyin: dze2 go3 / dze5 go3
- Guangdong Romanization: zé2 go3 / zé5 go3
- Sinological IPA (key): /t͡sɛː³⁵ kɔː³³/, /t͡sɛː¹³ kɔː³³/
Hakka
+ (Northern Sixian, incl. Miaoli)
- Pha̍k-fa-sṳ: liá-ke / yá-ke
- Hakka Romanization System: liaˋ ge / iaˋ ge
- Hagfa Pinyim: lia3 ge4 / ya3 ge4
- Sinological IPA: /li̯a³¹ ke⁵⁵/, /i̯a³¹ ke⁵⁵/
+ (Southern Sixian, incl. Neipu)
- Pha̍k-fa-sṳ: liá-ke
- Hakka Romanization System: liaˋ ge
- Hagfa Pinyim: lia3 ge4
- Sinological IPA: /li̯a³¹ ke⁵⁵/
+ (Meixian)
- Guangdong: zê3 gê4 / ê3 gê4 / gê3 gê4
- Sinological IPA: /t͡se³¹ ke⁵³/, /e³¹ ke⁵³/, /ke³¹ ke⁵³/
Note:
Meixian:
zê3 gê4 - literary;
ê3 gê4/gê3 gê4 - colloquial.
Southern Min
(Hokkien: Xiamen, Zhangzhou, General Taiwanese, Singapore)
Pe̍h-ōe-jī: chit-ê
Ti-lô: tsit-ê
Phofsit Daibuun: cit'ee
IPA (Kaohsiung): /t͡sit̚³²⁻⁴ e²³/
IPA (Xiamen, Taipei, Singapore): /t͡sit̚³²⁻⁴ e²⁴/
IPA (Zhangzhou): /t͡sit̚³²⁻⁵ e¹³/
(Hokkien: Quanzhou, Singapore)
Pe̍h-ōe-jī: chit-gê
Ti-lô: tsit-gê
Phofsit Daibuun: citgee
IPA (Quanzhou): /t͡sit̚⁵⁻²⁴ ɡe²⁴/
IPA (Singapore): /t͡sit̚³²⁻⁴ ɡe²⁴/
(Hokkien: Quanzhou, Jinjiang, Philippines)
Pe̍h-ōe-jī: chit-gé
Ti-lô: tsit-gé
Phofsit Daibuun: citgea
IPA (Quanzhou, Jinjiang, Philippines): /t͡sit̚⁵⁻²⁴ ɡe⁵⁵⁴/
(Hokkien: Jinjiang)
Pe̍h-ōe-jī: chiak-gê
Ti-lô: tsiak-gê
Phofsit Daibuun: ciakgee
IPA (Jinjiang): /t͡siak̚⁵⁻²⁴ ɡe²⁴/
(Hokkien: General Taiwanese)
Pe̍h-ōe-jī: chit-e
Ti-lô: tsit-e
Phofsit Daibuun: cit'ef
IPA (Taipei, Kaohsiung): /t͡sit̚³²⁻⁴ e⁴⁴/
(Hokkien: General Taiwanese)
Pe̍h-ōe-jī: chih-ge
Ti-lô: tsih-ge
Phofsit Daibuun: cihgef
IPA (Taipei): /t͡si(ʔ)³²⁻⁵³ ɡe⁴⁴/
IPA (Kaohsiung): /t͡si(ʔ)³²⁻⁴¹ ɡe⁴⁴/
(Hokkien: Quanzhou-like accent in Taiwan)
Pe̍h-ōe-jī: chih-gê
Ti-lô: tsih-gê
Phofsit Daibuun: cihgee
(Hokkien: Quanzhou-like accent in Taiwan)
Pe̍h-ōe-jī: chek-e
Ti-lô: tsik-e
Phofsit Daibuun: zek'ef
(Hokkien: Quanzhou-like accent in Taiwan)
Pe̍h-ōe-jī: chek-ê
Ti-lô: tsik-ê
Phofsit Daibuun: zek'ee
(Hokkien: Quanzhou-like accent in Taiwan)
Pe̍h-ōe-jī: chek-ge
Ti-lô: tsik-ge
Phofsit Daibuun: zekgef
(Hokkien: Quanzhou-like accent in Taiwan)
Pe̍h-ōe-jī: chek-gê
Ti-lô: tsik-gê
Phofsit Daibuun: zekgee
(Hokkien: General Taiwanese)
Pe̍h-ōe-jī: chit-le
Ti-lô: tsit-le
Phofsit Daibuun: citlef
IPA (Taipei, Kaohsiung): /t͡sit̚³²⁻⁴ le⁴⁴/
(Hokkien: General Taiwanese, Penang)
Pe̍h-ōe-jī: chit-lê
Ti-lô: tsit-lê
Phofsit Daibuun: citlee
IPA (Penang): /t͡sit̚³⁻⁴ le²³/
IPA (Taipei): /t͡sit̚³²⁻⁴ le²⁴/
IPA (Kaohsiung): /t͡sit̚³²⁻⁴ le²³/
(Hokkien: General Taiwanese)
Pe̍h-ōe-jī: chí-ê
Ti-lô: tsí-ê
Phofsit Daibuun: cyee
IPA (Kaohsiung): /t͡si⁴¹⁻⁴⁴ e²³/
IPA (Taipei): /t͡si⁵³⁻⁴⁴ e²⁴/
(Hokkien: Philippines)
Pe̍h-ōe-jī: chí-gé
Ti-lô: tsí-gé
Phofsit Daibuun: cy'gea
IPA (Philippines): /t͡si⁵⁵⁴⁻²⁴ ɡe⁵⁵⁴/
(Hokkien: Philippines)
Pe̍h-ōe-jī: chit-é
Ti-lô: tsit-é
Phofsit Daibuun: cit'ea
IPA (Philippines): /t͡sit̚⁵⁻²⁴ e⁵⁵⁴/
(Hokkien: Penang)
Pe̍h-ōe-jī: chih-lê
Ti-lô: tsih-lê
Phofsit Daibuun: cihlee
IPA (Penang): /t͡si(ʔ)³⁻⁴ le²³/
(Hokkien: Penang)
Pe̍h-ōe-jī: chí-lê
Ti-lô: tsí-lê
Phofsit Daibuun: cylee
IPA (Penang): /t͡si⁴⁴⁵⁻³³ le²³/
(Hokkien: Singapore)
Pe̍h-ōe-jī: che
Ti-lô: tse
Phofsit Daibuun: zef
IPA (Singapore): /t͡se⁴⁴/
(Hokkien: Philippines)
Pe̍h-ōe-jī: ché
Ti-lô: tsé
Phofsit Daibuun: zea
IPA (Philippines): /t͡se⁵⁵⁴/
Note: che/ché - contraction.
Pronoun
[edit]
這個
this; this one
: 這個酒店很衛生。 [MSC, trad.]
这个酒店很卫生。 [MSC, simp.]: Zhè ge jiǔdiàn hěn wèishēng. [Pinyin]
: This hotel is very hygienic.
: 他們為了解決這個問題,嘗試過各種方法。 [MSC, trad.]
他们为了解决这个问题,尝试过各种方法。 [MSC, simp.]: Tāmen wèile jiějué zhège wèntí, chángshì guò gèzhǒng fāngfǎ. [Pinyin]
: They have tried every method to solve this problem.
: 它在這個過程中起什麼作用? [MSC, trad.]
它在这个过程中起什么作用? [MSC, simp.]: Tā zài zhège guòchéng zhōng qǐ shénme zuòyòng? [Pinyin]
: What role does it play in this process?
: 這個公司在海外市場佔有一席之地。 [MSC, trad.]
这个公司在海外市场占有一席之地。 [MSC, simp.]: Zhège gōngsī zài hǎiwài shìchǎng zhànyǒu yīxízhīdì. [Pinyin]
: This company occupies a niche in the overseas market.
Usage notes
[edit]
Unlike English this, 這個/这个 (zhège) in Mandarin is not usually used for units of time such as 年 (nián), 天 (tiān) and 日 (rì). Rather, 今 (jīn) and 本 (běn) are used. Some common examples include 今晚 (jīnwǎn, “tonight”), 今天 (jīntiān, “today”) and 今年 (jīnnián, “this year”). However, 這個/这个 (zhège) can be used with 月 (yuè) - 這個月/这个月 (zhège yuè) - or more formally 本月 (běnyuè), as well as 星期 (xīngqī) - 這個星期/这个星期 (zhège xīngqī).
For Hokkien, Douglas (1873/1899) notes that dê/--dê is sometimes heard after final ⟨-t⟩, such as with this term, chit-ê / chit-le / chit-lê / chí-lê > chit-dê / chit-de / chid-dê / chí-dê, although modern Pe̍h-ōe-jī (POJ) does not use the letter ⟨D⟩.
Synonyms
[edit]
show ▼Dialectal synonyms of 這個 (“this; this one”) [map]
| Variety | Location | Words |
---
| Formal (Written Standard Chinese) | | 這個 |
| Northeastern Mandarin | Beijing | 這個 |
| Taiwan | 這個 |
| Harbin | 這個 |
| Singapore | 這個 |
| Jilu Mandarin | Jinan | 這個 |
| Jiaoliao Mandarin | Yantai (Muping) | 這個 |
| Central Plains Mandarin | Luoyang | 這個 |
| Xi'an | 這個 |
| Xining | 致個 |
| Xuzhou | 這個 |
| Lanyin Mandarin | Yinchuan | 這個 |
| Ürümqi | 這個 |
| Southwestern Mandarin | Chengdu | 這個 |
| Wuhan | 這個 |
| Guiyang | 這個 |
| Liuzhou | 這個 |
| Jianghuai Mandarin | Yangzhou | 這個 |
| Hefei | 這個 |
| Cantonese | Guangzhou | 呢個, 依個 |
| Hong Kong | 呢個, 依個 |
| Hong Kong (Kam Tin; Weitou) | 呢一個 |
| Macau | 呢個 |
| Guangzhou (Panyu) | 呢個 |
| Guangzhou (Huashan, Huadu) | 呢個 |
| Guangzhou (Zengcheng) | 呢個 |
| Foshan | 啲個 |
| Foshan (Shatou, Nanhai) | 離個 |
| Foshan (Shunde) | 呢個 |
| Foshan (Sanshui) | 哩個 |
| Foshan (Mingcheng, Gaoming) | 呢個 |
| Zhongshan (Shiqi) | 姑個 |
| Zhongshan (Huancheng) | 姑個 |
| Zhongshan (Nanlang) | 姑個 |
| Zhongshan (Xiaolan) | 呢個 |
| Zhongshan (Dongsheng, Xiaolan) | 呢個 |
| Zhongshan (Tanbei, Xiaolan) | 呢個 |
| Zhongshan (Henglan) | 呢個 |
| Zhongshan (Dongfeng) | 呢個 |
| Zhongshan (Nantou) | 呢個 |
| Zhongshan (Fusha) | 呢個 |
| Zhongshan (Shalang) | 呢個 |
| Zhongshan (Gangkou) | 呢個 |
| Zhongshan (Sanjiao) | 呢個 |
| Zhongshan (Guzhen) | 該個 |
| Zhongshan (Banfu) | 呢個 |
| Zhongshan (Tanzhou) | 呢個 |
| Zhuhai (Qianshan, Xiangzhou) | 呢個 |
| Zhuhai (Shangheng, Doumen; Tanka) | 呢個 |
| Zhuhai (Doumen) | 該個 |
| Jiangmen (Xinhui) | 該個 |
| Taishan | 該個 |
| Kaiping (Chikan) | 該個 |
| Heshan (Yayao) | 啟個 |
| Dongguan | 呢個 |
| Shenzhen (Shajing, Bao'an) | 呢個 |
| Qingyuan | 呢個 |
| Fogang | 呢個 |
| Yingde (Hanguang) | 呢個 |
| Yangshan | 利拉 |
| Shaoguan | 呢個 |
| Shaoguan (Qujiang) | 呢個 |
| Renhua | 呢個 |
| Lechang | 呢個 |
| Zhaoqing (Gaoyao) | 嗰個 |
| Sihui | 呢個 |
| Guangning | 嗰呢 |
| Deqing | 嗰個 |
| Huaiji | 嗰那 |
| Yunfu | 呢一個 |
| Luoding | 呢個 |
| Yunan (Pingtai) | 嗰個 |
| Yangjiang | 果個 |
| Xinyi | 己個隻 |
| Maoming (Xinpo) | 己隻 |
| Lianjiang | 嗰隻 |
| Kuala Lumpur (Guangfu) | 呢個 |
| Singapore (Guangfu) | 呢個, 依個 |
| Gan | Lichuan | 該個 |
| Pingxiang | 該隻 |
| Hakka | Meixian | 這隻, 這個 |
| Huizhou (Huicheng; Bendihua) | 呢隻 |
| Dongguan (Qingxi) | 惹隻 |
| Shenzhen (Shatoujiao) | 惹隻 |
| Zhongshan (Nanlang Heshui) | 該隻 |
| Guangzhou (Lütian, Conghua) | 底隻 |
| Miaoli (N. Sixian) | 這個, 這隻 |
| Pingtung (Neipu; S. Sixian) | 這個, 這隻 |
| Hsinchu County (Zhudong; Hailu) | 這個, 這隻 |
| Taichung (Dongshi; Dabu) | 這個, 這隻 |
| Hsinchu County (Qionglin; Raoping) | 這個, 這隻 |
| Yunlin (Lunbei; Zhao'an) | 這個, 這隻 |
| Jin | Taiyuan | 這, 這個 |
| Xinzhou | 這, 這個 |
| Northern Min | Jian'ou | 𱕇隻 |
| Eastern Min | Fuzhou | 只隻, 只一隻, 只回 |
| Southern Min | Xiamen | 這個 |
| Xiamen (Tong'an) | 這個 |
| Quanzhou | 這個, 種個 |
| Jinjiang | 這個, 種個 |
| Nan'an | 這個 |
| Shishi | 這個 |
| Hui'an | 這個 |
| Anxi | 這個 |
| Yongchun | 這個 |
| Dehua | 這個 |
| Zhangzhou | 這個 |
| Zhangzhou (Longhai) | 這個 |
| Zhangzhou (Changtai) | 這個 |
| Hua'an | 這個 |
| Nanjing | 這個 |
| Pinghe | 這個 |
| Zhangpu | 這個 |
| Yunxiao | 這個 |
| Zhao'an | 這個 |
| Dongshan | 這個 |
| Taipei | 這個 |
| Kaohsiung | 這個 |
| Yilan | 這個 |
| Changhua (Lukang) | 這個 |
| Taichung | 這個 |
| Taichung (Wuqi) | 這個 |
| Tainan | 這 |
| Taitung | 這個 |
| Hsinchu | 這個 |
| Kinmen | 這個 |
| Penghu (Magong) | 這個 |
| Singapore (Hokkien) | 這個 |
| Manila (Hokkien) | 這個, 種個 |
| Zhangping | 許個 |
| Datian | 這個 |
| Chaozhou | 只個 |
| Shantou | 只個 |
| Shantou (Chenghai) | 只個 |
| Shantou (Chaoyang) | 只個 |
| Jieyang | 只個 |
| Johor Bahru (Teochew) | 只個 |
| Singapore (Teochew) | 只個 |
| Batam (Teochew) | 只個 |
| Wenchang | 這個, 這枚 |
| Qionghai | 這枚 |
| Singapore (Hainanese) | 這個, 這枚 |
| Zhongshan Min | Zhongshan (Longdu, Shaxi) | 指個 |
| Wu | Shanghai | 搿個, 迭個 old-style |
| Shanghai (Fengxian) | 格當 |
| Shanghai (Jinshan) | 格當, 格個 |
| Suzhou | 該格, 搿格 |
| Hangzhou | 接個, 格個 |
| Ningbo | 該個 |
| Wenzhou | 該個, 個個 |
| Xiang | Changsha | 咯個, 咯隻 |
| Shuangfeng | 咯隻 |
Descendants
[edit]
Wutunhua: je-ge
Interjection
[edit]
這個
uh; um; er; you know
Synonyms
[edit]
那個/那个 (nèige)
References
[edit]
^ Douglas, Carstairs (1873) “dê”, in Chinese-English Dictionary of the Vernacular or Spoken Language of Amoy, With the Principal Variations of the Chang-chew and Chin-chew Dialects. (overall work in Hokkien and English), London: Trübner & Co., page 99; New Edition, With Corrections by the Author., Thomas Barclay, Lîm Iàn-sîn 林燕臣, London: Publishing Office of the Presbyterian Church of England, 1899, page 99
Retrieved from "
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13792 | https://www.cuemath.com/numbers/identity-property/ | LearnPracticeDownload
Identity Property
Identity property is applied to a group of numbers in the form of sets. The identity of these numbers remains the same as 1 and 0 even when the numbers are added, subtracted, multiplied, and divided. Among the arithmetic operations, addition and multiplication are used more often. Let us learn more about the identity property, types, and solve a few examples.
| | |
--- |
| 1. | Identity Property Definition |
| 2. | Additive Identity Property |
| 3. | Multiplicative Identity Property |
| 4. | Additive Vs Multiplicative |
| 5. | FAQs on Identity Property |
Identity Property Definition
Identity property is defined as the property where if any arithmetic operations are used to combine an identity with a number (n), the end result will be n. Any number can be considered as an identity when added, subtracted, multiplied, or divided with another number such as n, allows n to remain the same. This identity or end result will always be a 0 or 1 depending on the arithmetic operation used. In the case of addition and subtraction, the result is 0 i.e. when 0 is added or subtracted from n, then the value of n remains the same. Whereas in multiplication and division, the identity is 1 i.e. when 1 is multiplied or divided with n, the value of n remains the same. Look at the image below to understand the concept better.
Identity is Always 0 and 1
As we already learned that the identity will always be 0 and 1 depending on the arithmetic operation used. The reason for this is as follows:
In mathematics, addition is considered as the process of adding any one number to another number. According to the identity property, when a number is added to 0, the result will be the same. For example, 25 + 0 = 25.
Subtraction is the process of taking away one number from the given whole number. Hence, when 0 is subtracted from a number, the result is the same. For example, 16 - 0 = 16.
Multiplication is the process of repeated addition i.e. the number is repeatedly added by the number of times stated by the other number. For example, 6 × 5, 6 is repeatedly being added 5 times. According to the identity property of multiplication, when a number is multiplied by the number 1, the product will be the number itself, 6 × 1 = 6.
Division is the process of separating a number into parts. According to the identity property, when a number is divided by 1, the result will be the number itself. For example, 4 ÷ 1 = 4.
Additive Identity Property
The additive identity property states that when a number is added to 0 then the result obtained will be the same as the number. Here, 0 is known as the identity element. This property is applicable for real numbers, complex numbers, integers, rational numbers, whole numbers, and so on. For example, 23 + 0 = 23.
Therefore, the additive identity formula is a + 0 = a (a is any real number)
Multiplicative Identity Property
The multiplicative identity property states that when a number is multiplied by 1, the product will be the number itself. Here, 1 is considered as the identity element. This property can be applied to real numbers, complex numbers, and so on. The multiplicative property is applicable for any non-zero rational number p/q and in the case of integers. For example, 87 × 87 = 87, where 87 is the number on which we applied the multiplicative identity.
Note: The multiplicative identity is not applied when any number is multiplied by -1 because the result will not be the same number. For example, 29 × -1 = -29.
The multiplicative identity formula is expressed as a × 1 = a (a is any real number)
Additive Vs Multiplicative
The below table shows the difference between additive identity and multiplicative identity.
| Additive Identity Property | Multiplicative Identity Property |
--- |
| Used in addition | Used in multiplication |
| The additive identity element is 0 | The multiplicative identity element is 1 |
| Formula: a + 0 = a | Formula: a × 1 = a |
☛ Related Articles
Listed below are a few articles related to identity property.
Zero Property of Multiplication
Associative Property of Addition
Additive Identity vs Multiplicative Identity
Associative Property of Multiplication
Read More
Download FREE Study Materials
Identity Property Worksheets
Identity Property Worksheets
Worksheets on Identity Property
Identity Property Examples
Example 1: Select the equation that satisfies the multiplicative identity property.
a) 5/7 × 1 = 5/7
b) 9 × 0 = 0
Solution: According to the multiplicative identity property when we multiply any rational number by 1 the result will be the same rational number.
a) 5/7 × 1 = 5/7, this equation satisfies the property because the result is the same number that is 5/7 and the multiplicative identity element is 1 in this case.
b) 9 × 0 = 0, this equation does not satisfy the property as the result is not the same number 9. The result is 0 and the multiplicative identity element is not 1 in this case.
Therefore, (a) satisfies the multiplicative identity property.
2. Example 2: A flock of 13 penguins sat together. No more penguins joined the lot. Find the number of penguins sitting together using additive identity.
Solution: As the condition says no penguins joined the flock of 13 penguins, hence the number of penguins remained the same.
= 13 penguins + 0 penguins = 13 penguins
3. Example 3: State true or false.
a) The additive identity property is also known as the zero property of addition.
b) According to the multiplicative identity property, multiplying 0 to any number, results in the number itself.
Solution:
a) True, the additive identity property is also known as the zero property of addition.
b) False, according to the multiplicative identity property, multiplying 1 to any number, results in the number itself.
View Answer >
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Practice Questions on Identity Property
Check Answer >
FAQs on Identity Property
What is Identity Property?
Identity property states that when any number is combined with an identity either 0 or 1, the end result will be the number itself. The property is applicable while using the four main arithmetic operations - addition, multiplication, subtraction, and division.
What is the Additive Identity Property in Math?
The additive identity property states that when a number is added to zero, it results in the number itself. For example, if 8 is added to 0, the sum is the number itself. 8 + 0 = 8. Here, zero is known as the identity element which keeps the identity of the number.
What does the Multiplicative Identity Property Mean in Math?
According to the multiplicative identity property, if a number is multiplied by 1, it results in the number itself. For example, if 12 is multiplied by 1, the product is the number itself (12 × 1 = 12). Here, one is known as the identity element which keeps the identity of the number.
What is the Difference Between Additive Identity and Multiplicative Identity?
The following points show the differences between Additive Identity and Multiplicative Identity.
The additive identity of numbers is used for the addition operation, whereas, the multiplicative identity is used for a multiplication operation.
0 is the identity element in the additive identity (a + 0 = a), whereas, 1 is the identity element in the multiplicative identity (a × 1 = a).
25 + 0 = 25 is the example of the additive identity property and 25 × 1 = 25 is the example of the multiplicative identity property.
What is the Identity Element in Additive Identity and Multiplicative Identity?
0 is the identity element in additive identity.
1 is the identity element in the multiplicative identity.
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13794 | https://www.youtube.com/watch?v=d4CdMoucVfs | Calculus - Definite Integrals
The Organic Chemistry Tutor
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Posted: 21 Dec 2019
This calculus video tutorial explains how to evaluate a definite integral. It also explains the difference between definite integrals and indefinite integrals.
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in this video we're going to talk about how to evaluate definite integrals now before we begin take a minute to hit that subscribe button and don't forget to turn on the notification bell so let's talk about the difference between a definite integral and an indefinite integral a definite integral has a lower limit of integration in this case a and b the upper limit of integration an indefinite integral does not have that so this is an indefinite integral and this here is a definite integral the antiderivative of a function f of x is capital f to evaluate the definite integral once you find the antiderivative you need to plug in the limits of integration and so the the value of the definite integral is going to be f of b minus f of a but let me show you the process by which we can evaluate a definite integral so let's start with this example we need to find the antiderivative of each term in this expression what is the antiderivative of 8x cubed well let me give you a review first the antiderivative of a variable raised to a constant is going to equal that variable raised to the constant plus 1 divided by n plus 1. now for indefinite integrals we would have the constant of integration c but when dealing with definite integrals you don't need to worry about c so as an example let's say if we want to find the anti-derivative of x to the fifth all we need to do is add one to the exponent it's going to be x to the sixth and then divide by that result so let's say if we want to determine the antiderivative of 4 x to the 7th so we have a constant times x to the 7th first rewrite the constant and then find the antiderivative of x to the seventh so add one to seven that's eight and then divide by that number and after that you can reduce it so eight is four times two we can cancel the four and so the answer is going to be x to the eight over two so that's how you could find the antiderivative of monomials now let's continue on with this example so to find the anti-derivative of 8x cubed first we're going to rewrite the constant eight then we're going to add one to the exponent three plus one is four and then we're going to divide by 4. now let's repeat this process for the next one so the antiderivative of 3x squared is going to be the constant 3 times x raised to the third power divided by three now what about the anti-derivative of six times x if you don't see a number it's always a one this is six times x to the first power so just like before we're going to rewrite the constant 6 and then the variable add 1 to the exponent 1 plus 1 is 2 and then divide by that result now as was mentioned before because we're dealing with a definite integral we don't need to write the constant c here but we do need to write our limits of integration so now let's simplify this expression 8 divided by four is two so we have two times x to the fourth three divided by three is one so that cancels so we have one x cubed which we can write as just x cubed six divided by two is three so then this is going to be plus three x squared evaluated from two to three now this is going to equal f of three minus f of two and keep in mind this expression here represents lowercase f of x and this expression is the antiderivative which represents capital f of x so we're going to plug in 3 and 2 into capital f of x so in this in these brackets we're going to put f of 3 and here this is going to be f of 2. so let's plug in 3 into this expression so it's going to be 2 times 3 raised to the fourth power plus 3 raised to the third power plus 3 times 3 squared now let's substitute x with two in the second set of brackets so we have two times two to the fourth power plus two to the third plus 3 times 2 squared so this is f of 3 and this here is f of 2. so at this point we just need to do the math three to the fourth power is eighty-one and eighty-one times two is one sixty-two three to the third is twenty-seven three squared is nine times three that's twenty-seven as well two to the fourth power if you multiply four twos you're going to get sixteen and sixteen times two is thirty two two to the third power is eight two squared is four times three that's twelve 27 plus 27 we no longer need the brackets anymore 27 plus 27 is 54. 32 plus 8 is 40 and 40 plus 12 is 52 so this is minus 52 54 minus 52 is 2 162 plus 2 is 164. so this is the value of the definite integral now this value here 164 represents the area under the curve that is between the curve represented by that function and the x-axis between the x-values 2 and 3. so that's what the definite integral can do it can help you calculate the area under the curve but that's the topic for another discussion for those of you who want more examples on evaluating definite integrals check out the description section of this video i'm going to post some links there if you wish to find more examples even harder examples including square roots and other stuff so feel free to take a look at that and if you haven't done so already don't forget to subscribe to this channel and click on that notification bell thanks again for watching |
13795 | https://engineerstutor.com/2018/10/30/inverse-z-transform-power-series-expansion-method/ | Inverse z Transform | Power Series Expansion Method - EngineersTutor
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College Algebra 2e
Practice Test
College Algebra 2ePractice Test
Contents Contents
Highlights
Table of contents
Preface
1 Prerequisites
2 Equations and Inequalities
3 Functions
4 Linear Functions
5 Polynomial and Rational Functions
6 Exponential and Logarithmic Functions
7 Systems of Equations and Inequalities
8 Analytic Geometry
Introduction to Analytic Geometry
8.1 The Ellipse
8.2 The Hyperbola
8.3 The Parabola
8.4 Rotation of Axes
8.5 Conic Sections in Polar Coordinates
Chapter Review
Exercises
Review Exercises
Practice Test
9 Sequences, Probability, and Counting Theory
Answer Key
Index
Search for key terms or text.
Close
Practice Test
For the following exercises, write the equation in standard form and state the center, vertices, and foci.
1.
x 2 9+y 2 4=1 x 2 9+y 2 4=1 x 2 9+y 2 4=1
9 y 2+16 x 2−36 y+32 x−92=0 9 y 2+16 x 2−36 y+32 x−92=0 9 y 2+16 x 2−36 y+32 x−92=0
For the following exercises, sketch the graph, identifying the center, vertices, and foci.
3.
(x−3)2 64+(y−2)2 36=1(x−3)2 64+(y−2)2 36=1(x−3)2 64+(y−2)2 36=1
2 x 2+y 2+8 x−6 y−7=0 2 x 2+y 2+8 x−6 y−7=0 2 x 2+y 2+8 x−6 y−7=0
5.
Write the standard form equation of an ellipse with a center at (1,2),(1,2),(1,2), vertex at (7,2),(7,2),(7,2), and focus at (4,2).(4,2).(4,2).
A whispering gallery is to be constructed with a length of 150 feet. If the foci are to be located 20 feet away from the wall, how high should the ceiling be?
For the following exercises, write the equation of the hyperbola in standard form, and give the center, vertices, foci, and asymptotes.
7.
x 2 49−y 2 81=1 x 2 49−y 2 81=1 x 2 49−y 2 81=1
16 y 2−9 x 2+128 y+112=0 16 y 2−9 x 2+128 y+112=0 16 y 2−9 x 2+128 y+112=0
For the following exercises, graph the hyperbola, noting its center, vertices, and foci. State the equations of the asymptotes.
9.
(x−3)2 25−(y+3)2 1=1(x−3)2 25−(y+3)2 1=1(x−3)2 25−(y+3)2 1=1
y 2−x 2+4 y−4 x−18=0 y 2−x 2+4 y−4 x−18=0 y 2−x 2+4 y−4 x−18=0
11.
Write the standard form equation of a hyperbola with foci at (1,0)(1,0)(1,0) and (1,6),(1,6),(1,6), and a vertex at (1,2).(1,2).(1,2).
For the following exercises, write the equation of the parabola in standard form, and give the vertex, focus, and equation of the directrix.
y 2+10 x=0 y 2+10 x=0 y 2+10 x=0
13.
3 x 2−12 x−y+11=0 3 x 2−12 x−y+11=0 3 x 2−12 x−y+11=0
For the following exercises, graph the parabola, labeling the vertex, focus, and directrix.
(x−1)2=−4(y+3)(x−1)2=−4(y+3)(x−1)2=−4(y+3)
15.
y 2+8 x−8 y+40=0 y 2+8 x−8 y+40=0 y 2+8 x−8 y+40=0
Write the equation of a parabola with a focus at (2,3)(2,3)(2,3) and directrix y=−1.y=−1.y=−1.
17.
A searchlight is shaped like a paraboloid of revolution. If the light source is located 1.5 feet from the base along the axis of symmetry, and the depth of the searchlight is 3 feet, what should the width of the opening be?
For the following exercises, determine which conic section is represented by the given equation, and then determine the angle θ θ θ that will eliminate the x y x y x y term.
3 x 2−2 x y+3 y 2=4 3 x 2−2 x y+3 y 2=4 3 x 2−2 x y+3 y 2=4
19.
x 2+4 x y+4 y 2+6 x−8 y=0 x 2+4 x y+4 y 2+6 x−8 y=0 x 2+4 x y+4 y 2+6 x−8 y=0
For the following exercises, rewrite in the x′y′x′y′x′y′ system without the x′y′x′y′x′y′ term, and graph the rotated graph.
11 x 2+10 3–√x y+y 2=4 11 x 2+10 3 x y+y 2=4 11 x 2+10 3 x y+y 2=4
21.
16 x 2+24 x y+9 y 2−125 x=0 16 x 2+24 x y+9 y 2−125 x=0 16 x 2+24 x y+9 y 2−125 x=0
For the following exercises, identify the conic with focus at the origin, and then give the directrix and eccentricity.
r=3 2−sin θ r=3 2−sin θ r=3 2−sin θ
23.
r=5 4+6 cos θ r=5 4+6 cos θ r=5 4+6 cos θ
For the following exercises, graph the given conic section. If it is a parabola, label vertex, focus, and directrix. If it is an ellipse or a hyperbola, label vertices and foci.
r=12 4−8 sin θ r=12 4−8 sin θ r=12 4−8 sin θ
25.
r=2 4+4 sin θ r=2 4+4 sin θ r=2 4+4 sin θ
Find a polar equation of the conic with focus at the origin, eccentricity of e=2,e=2,e=2, and directrix: x=3.x=3.x=3.
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Authors: Jay Abramson
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Book title: College Algebra 2e
Publication date: Dec 21, 2021
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13798 | https://arxiv.org/pdf/1912.10462 | Published Time: Mon, 23 Jan 2023 03:22:13 GMT
arXiv:1912.10462v2 [math.NT] 12 Jul 2020
Lattice points in d-dimensional spherical segments
Martin Ortiz Ramirez
Abstract
We study lattice points in d-dimensional spheres, and count their number in thin spherical segments. We found an upper bound depending only on the radius of the sphere and opening angle of the segment. To obtain this bound we slice the segment by hyperplanes of rational direction, and then cover an arbitrary segment with one having rational direction. Diophantine approximation can be used to obtain the best rational direction possible.
Keywords: discrete geometry, lattice points on spheres, Diophantine approximation.
MSC(2010): 11P21, 11K60.
1 Introduction
The topic of lattice points in d-spheres has been extensively studied for centuries, starting with the two dimensional case of Gauss’s circle problem, which asks for the number of lattice points
N(r) inside a circle of radius r, with r2 = n an integer. Gauss already knew that by geometric considerations N(r) is close to the area of the circle πr 2, and the problem transitioned to the study of the asymptotic behaviour of E(r) = |N(r) − πr 2|. One natural question asks for the the infimum
μ of the α such that E(r) = O(rα). Gauss proved that μ ≤ 1 by considering the length of the circumference, Sierpinski improved it to μ ≤ 2/3, and Van der Corput showed that μ < 2/3 [8, p.20-22]. They also conjectured that μ = 1 /2. The best upper bound to date is μ ≤ 517 /824 by Bourgain and Watt . As for a lower bound, Hardy and Landau independently proved that
μ ≥ 1/2. Gauss’s circle problem can be interpreted by considering the sum of all the r2(m) for m ≤ n,where r2(m) is the number of integral solutions to x2 + y2 = m. Unlike Gauss’s problem, where obtaining the correct order of magnitude is easy, it is more difficult for r2(m), since it vanishes for arbitrarily large m. Relating r2(m) to the number of divisors of m [9, Ch 16.10] yields r2(m) ≪ mǫ
for any ǫ > 0 1. Consider the number of lattice points on a short circular arc. Cilleruelo and Cordoba proved that on a circle of radius R, an arc of length no greater than √2R1/2−1/(4 ⌊l/ 2⌋+2)
contains at most l lattice points. Cilleruelo and Granville proved that on an arc of length less
1Here and in what follows f(m)≪g(m) means that f(m)/g (m) is bounded as mtends to infinity.
1than (40 + 40
3
√10) 1/3R1/3 there are at most 3 lattice points, an improvement on a result by Jarnik , that in an arc of length less than or equal to than R1/3 there are at most two lattice points. Moving to 3 dimensions, S(R) denoting the number of lattice points inside a sphere of radius R,we have the analogue of Gauss’s circle problem. By the same geometric considerations S(R) is of the order of 4
3
πR 3, the volume of the sphere. Hence we are interested in E(R) = | 4
3
πR 3 − S(R)|.Szego proved that E(R) = o(R(logR )1/2) is not true 2. Regarding upper bounds of the form
E(R) = O(Rθ), Heath-Brown proved that it holds true with θ = 21 /16. As in the two dimensional case, the number r3(R2) of lattice points on the sphere of radius R,
R2 = n still vanishes for arbitrarily large n, of the form 4 a(8 k + 7) (Legendre’s theorem [8, Ch. 4]), although Walfisz proved that r3(R2) ≫ R log log R holds for infinitely many n. We also have the upper bound r3(R2) ≪ R1+ ǫ for any ǫ > 0 and the lower bound r3(R2) ≫ R1−ǫ for n 6 = 0 , 4, 7 mod 8 [3, 1] i.e. when there are primitive solutions to x2 + y2 + z2 = n (solutions where the greatest common factor of x, y and z is 1). The distribution of lattice points on the sphere has also been a topic of study. In the 1950s Linnik proved that as n tends to infinity amongst n square-free and
n ≡ ± 1 mod 5 3, the projections of the lattice points onto the unit sphere becomes equidistributed. Duke proved the result removing the constraint n ≡ ± 1 mod 5 after thirty years. A three dimensional analogue of an arc is a cap , the surface obtained by cutting the sphere with a plane, or equivalently, the intersection of RS 2 and sphere of radius λ centered at a point of the sphere. Bourgain and Rudnick proved using a theorem by Jarnik that the maximal number
F3(R, λ ) of lattice points in a cap of radius λ satisfies
F3(R, λ ) ≪ Rǫ
(
1 + λ2
R1/2
)
for all ǫ > 0.
This is only nontrivial for small caps having λ ≪ R1/2, because other methods already give
F3(R, λ ) ≪ Rǫ(1 + λ) [2, Lemma 2.2]. For bigger caps Bourgain and Rudnick also proved [2, Proposition 1.3]
F3(R, λ ) ≪ Rǫ
(
1 + λ
( λ
R
)η)
for any 0 < η < 1
15 .
If we slice the sphere by two parallel planes we obtain a spherical segment . Maffucci [13, Propo-sition 6.2] gave a bound (see the paragraph before Theorem 1.3) for the number of lattice points in segments given their opening angle, a theorem that we generalise in this paper. The four dimensional case still shows an erratic behaviour, for instance r4(2 n) = 24 for all n. For
d greater than 4 the behaviour of rd is much nicer, as the order of magnitude of rd is known to be
rd(R2) ≈ Rd−2, via the circle method [8, Chapters 10-12]. By similar methods Bourgain-Rudnick [2, Appendix A] also proved bounds for the number of lattice points in d-dimensional spherical caps:
Fd(R, λ ) ≪ Rǫ
( λd−1
R + λd−3
)
for d ≥ 5.
2Here and elsewhere f(R)=o(g(R)) means that f(R)/g(R) tends to 0 as Rtends to infinity.
3In fact, Linnik aslo proved the statement replacing the constraint mod 5 by a similar one mod p, for any fixed prime p.
2The proof for d ≥ 5 is considerably different from its three dimensional analogue. The former uses analytic tools such as the circle method. The latter is more geometric in nature, and uses Diophantine approximation, in a similar way to this paper. While interesting in its own right, lattice points problems have found applications in other fields. Recently they have been applied to the study of arithmetic waves, with results by Oravecz-Rudnick-Wigman , Krishnapur-Kurlberg-Wigman , Rossi-Wigman and Benatar-Maffucci among others. Before stating the main results, we first need to define some key concepts. Let Sd−1 denote the unit ( d − 1)-sphere in Rd, and B(α, r ) the closed ball of radius r around α ∈ Rd. For a vector
β ∈ Rd, |β| will denote its euclidean norm.
Definition 1.1. A spherical cap T ⊆ RS d−1 of direction β ∈ Rd, |β| = 1 , and radius r is defined as
T = RS d−1 ∩ B(Rβ, r ).
We then define a spherical segment S ⊆ RS d−1 as S = T1 \ T2, where each Ti is a spherical cap of direction β and radius ri, r1 > r 2. We also say that the direction of any of the Ti is the direction of
S.
We will define RS d−1 ∩ { x ∈ Rd : |x − Rβ | = r} as the base of T , we will see in the next section that it is a ( d − 2)-sphere in a hyperplane orthogonal to β. In a ( d − 1)-sphere with centre O and radius r, two points P1, P 2 are said to be antipodal if P1 = O + v, P2 = O − v, where v has norm r.Hence the following definition.
Definition 1.2. The opening angle of a spherical cap is ∠P OQ , where O is the origin and P ,Q
are antipodal points of the (d − 2) -sphere that is the base. The opening angle θ of a segment S is
θ1 − θ2, where θi is the opening angle of Ti.
We define ψ(R, θ ) to be the maximal number of lattice points in segments T ⊆ RS d−1 of opening angle θ. Our method will involve slicing the segment by hyperplanes. In this context, κd(R) 4 is the maximal number of integer points in the intersection of RS d−1 and a rational hyperplane , that is, a hyperplane that can be defined by a the equation of the form a1x1 + a2x2 + . . . a dxd = a, for ai ∈ Z.We are now in position to state the main results of the paper. The first one is a generalisation of the following result by Maffucci [13, Proposition 6.2] for the three dimensional case:
ψ ≪ κ3(R)(1 + Rθ 1/3) as θ → 0 along with R → ∞ .
In all future statements Ψ( R, θ ) ≪ f (R, θ ) as θ → 0 will mean that θ tends to 0 together with
R → ∞ . Furthermore, unless otherwise stated, the result is understood to be uniform, that is, the implied constant does not depend on R or θ.
4To the best of my knowledge, an upper bound for κdwhen d≥4 has not been explicitly studied in the literature, but a bound like κd(R)≪Rd−3+ ǫwould be expected to be amenable to the circle method for d≥5, since the intersection of a hyperplane and a ( d−1)-sphere is a ( d−2)-sphere.
3Theorem 1.3. Let ψ(R, θ ) be the maximal number of lattice points on spherical segments in RS d−1
with opening angle θ. Then
ψ ≪ κd(R)(1 + Rθ 1/d )
as θ → 0.
We can improve Theorem 1.3 if we are given more information about the direction of the segment.
Definition 1.4. For β ∈ Rd, β is said to have s rational quotients if there exists a non-zero k ∈ R
such that kβ has exactly s + 1 rational coordinates, and this number of rational coordinates is maximal amongst all the non-zero k.
We then fix some direction β and ask for the number of lattice points in segments of radius R
and opening angle θ. The case s = 0 is already covered in Theorem 1.3.
Theorem 1.5. Let S ⊆ RS d−1 be a spherical segment of opening angle θ and direction β having s
rational quotients, 1 ≤ s ≤ d − 1. Let ψ(R, θ, β ) be the number of lattice points in S. Then
ψ ≪β κd(R)
(
1 + Rθ 1
d−s
)
as θ → 0. 5
Next, we prove some general geometric facts about segments that will be useful later on. On Section 3 we will state all the relevant lemmas and prove Theorem 1.3. Section 4 is devoted to the proof of such lemmas and we prove Theorem 1.5 on Section 5.
2 Some geometric considerations
In this section we introduce a few geometric features of spherical segments that will be used later on. Recall Definition 1.1, we now prove that the base of each Ti is indeed a ( d − 2)-sphere. Let v · w
denote the dot product of two vectors in Rd. The intersection of {| x − Rβ | = ri} and {| x| = R}
lies in the plane β · x = R − r2
i
2R
:= λi since
|x − Rβ |2 = R2 + |x|2 − 2Rβ · x = r2
i
holds on the intersection. Then
|λiβ − x|2 = λ2
i
R2 − 2λiβ · x = R2 − λ2
i
= r2
i
− r4
i
4R2 ,
5Here ≪βmeans that the implied constant depends on β.
4so that S has two bases B1 and B2 that are ( d − 2)-spheres lying on the hyperplanes β · x = λi, of radii
k2
i
:= r2
i
− r4
i
4R2 .
In Definition 1.2, P OQ is an isosceles triangle of side lengths R and 2 ki, so θ is well-defined. Because P and Q are antipodal points, O, P, Q and Rβ are coplanar, hence by basic geometry we have that
ri = 2 R sin( θi/4) . (2.1) Another parameter that will be useful is the height of the segment, the distance between the two bases B1 and B2
h =
∣∣∣∣
βλ 1
|β|2 − βλ 2
|β|2
∣∣∣∣ = 1
|β||λ1 − λ2| = 1
2R (r21 − r22). (2.2)
3 Proof of Theorem 1.3
We will first state the necessary lemmas, that will be proved in the next section.
Lemma 3.1. Let S ⊆ RS d−1 be a spherical segment of direction b
|b|
, with b ∈ Zd and height
0 ≤ h ≤ 2R. Then
ψ ≤ κd(R)(1 + |b|h).
Lemma 3.1, although simple, provides the best upper bound using slicing methods for a segment of rational direction; we will use it to prove the next lemma. Lemma 3.1 will be proved in the next section.
Lemma 3.2. Let S ⊆ RS d−1 be a spherical segment with direction β and opening angle θ. For any
a ∈ Zd, the maximal number of lattice points lying on S satisfies
ψ ≪ κd(R)(1 + R|a|(θ + φ)) as θ → 0
where 0 ≤ φ ≤ π is the angle between β and a. The implied constant is absolute.
We will prove Lemma 3.2 in the next section. We will try to find an a ∈ Zd that optimises |a|
and φ simultaneously using Diophantine approximation. The main ingredient is the following.
Theorem 3.3. (Dirichlet’s theorem) [16, Theorem 1A, p.27]. Let ξi ∈ R for i = 1 , 2 . . . d , and H ∈ N. Then there exist pi ∈ Z and 1 ≤ q ≤ Hd such that
∣∣∣∣ξi − pi
q
∣∣∣∣ ≤ 1
qH ∀ i = 1 , 2 . . . d.
Note that
∣∣∣ a
|a|
− β
∣∣∣ = 2 sin ( φ/ 2) ∼ φ as φ → 0. The next lemma will provide a good simultaneous bound for |a| and φ.5Lemma 3.4. For all β ∈ Rd, |β| = 1 , and integers H ≥ 1, there exists some a ∈ Zd such that
|a| ≪ Hd−1 and
∣∣∣∣
a
|a| − β
∣∣∣∣ ≪ 1
|a|H ,
where ≪ is understood with respect to H, and the implied constants are absolute.
We will prove Lemma 3.4 in the next section.
Proof of Theorem 1.3 assuming all the lemmas. We follow the proof of the 3-dimensional case by Maffucci [13, Prop. 6.2]. Given a segment with direction β we choose a ∈ Zd satisfying Lemma 3.4, for some H to be determined. Then one has
∣∣∣∣β − a
|a|
∣∣∣∣ = 2 sin( φ/ 2) ∼ φ as φ → 0.
Later we will choose H so that θ → 0 implies φ → 0, in which case by Lemma 3.2 gives
ψ ≪ κd(R)
(
1 + R|a|θ + R|a|
∣∣∣∣β − a
|a|
∣∣∣∣)
as θ → 0. (3.1) It follows from Lemma 3.4 that
|a|θ + |a|
∣∣∣∣β − a
|a|
∣∣∣∣ ≪ θH d−1 + 1
H . (3.2) Setting H = ⌈θ−1/d ⌉ in (3.2), so that H = O(θ−1/d ), we obtain
|a|θ + |a|
∣∣∣∣β − a
|a|
∣∣∣∣ ≪ θ1/d .
Therefore, since the implied constants in (3.1) and (3.2) are absolute, if follows that
ψ ≪ κd(R)(1 + Rθ 1/d ) as θ → 0.
To complete the proof we have that for our choice of H,2 sin( φ/ 2) =
∣∣∣∣β − a
|a|
∣∣∣∣ ≪ 1
|a|H ≤ 1
H ∼ θ1/d → 0 as θ → 0.
Therefore φ → 0 as θ → 0, because 0 ≤ φ ≤ π.
4 Proofs of the Lemmas
Proof of Lemma 3.1. Since b ∈ Zd, all lattice points in the segment lie on a (rational) hyperplane of the form b · x = n ∈ Z. The distance between planes is 6∣∣∣∣
nb
|b|2 − (n + 1) b
|b|2
∣∣∣∣ = 1
|b|
and there are at most 1 +
⌊ h
1/|b|
⌋
≤ 1 + |b|h
hyperplanes intersecting the segment, each of them with at most κd(R) lattice points in them; thus the lemma follows.
A slightly different proof of Lemma 3.1 when d = 3 can be found in [13, Proposition 6.3].
Proof of Lemma 3.2. Let S = T1 \ T2 with respective radii r1 and r2. We will construct S′ = T ′
1
\ T ′
2
having direction proportional to a ∈ Zd, such that S ⊆ S′, then ψ can be bounded above using Lemma 3.1 on S′. For this it will suffice to construct S′ so that T ′
2
⊆ T2 and T1 ⊆ T ′
1
. Let r′
1
, r ′
2
be the radii of T ′
1
and T ′
2
respectively. We claim that taking
r′
2
= r2 − 2R sin( φ/ 2) and
r′21 = r21 + 4 Rr 1 sin( φ/ 2) + 4 R2 sin( φ/ 2) 2
satisfies the conditions. A point inside T ′
2
is of the form R a
|a|
v with |v| ≤ r′
2
, hence
∣∣∣∣Rβ −
(
R a
|a| + v
)∣ ∣∣∣
2
= R2
∣∣∣∣β − a
|a|
∣∣∣∣
2
|v|2 + 2 Rv ·
(
β − a
|a|
)
≤ 4R2 sin( φ/ 2) 2 + |v|2 + 2 R|v|
∣∣∣∣β − a
|a|
∣∣∣∣
≤ 4R2 sin( φ/ 2) 2 + r′22 + 4 R sin( φ/ 2) r′
2
= r22,
where we have used that
∣∣∣β − a
|a|
∣∣∣ = 2 sin( φ/ 2). This shows that T ′
2
⊆ T2. Now let Rβ + v with
|v|≤ r1 be a point of T1. We have
∣∣∣∣(Rβ + v) − R a
|a|
∣∣∣∣
2
≤ 4R2 sin( φ/ 2) 2 + |v|2 + 2 R|v|
∣∣∣∣β − a
|a|
∣∣∣∣
≤ 4R2 sin( φ/ 2) 2 + r21 + 4 Rr 1 sin( φ/ 2) = r′21 ,
so that T1 ⊆ T ′
1
as desired. Now, according to (2.2), the height of S′ is
h′ = 1
2R (r′21 − r′22 ) = 1
2R
(r21 − r22 + 4 R sin( φ/ 2)( r1 + r2)) .
Recall that ri = 2 R sin( θi/4) (2.1), and θ = θ1 − θ2. Therefore,
h′ = 2 R(sin( θ1/4) 2 − sin( θ2/4) 2) + 4 R sin( φ/ 2)(sin( θ1/4) + sin( θ2/4)) . (4.1) 7Now, θ1/4 = θ2/4 + θ/ 4, so that using the Taylor expansion of sin( x) yields sin( θ1/4) = sin( θ2/4) cos( θ/ 4) + cos( θ2/4) sin( θ/ 4) θ
4 + O(θ2),
implying sin( θ1/4) 2 = sin( θ2/4) 2 + sin( θ2/2) θ
4 + O(θ2). (4.2) It follows from (4.1) and (4.2) that
h′ = 2 R
(
sin( θ2/2) θ
4 + 2 sin( φ/ 2) (sin( θ1/4) + sin( θ2/4)) + O(θ2)
)
. (4.3) We have that sin( φ/ 2) ≤ φ/ 2, and all the other sines in (4.3) are bounded above by 1, thus
h′ ≪ R(θ + φ)as R tends to infinity and θ tends to 0, and the implied constant is absolute. By Lemma 3.1, the number of lattice points in S′ is no greater than
κd(R)(1 + |a|h′) ≪ κd(R)(1 + R|a|(θ + φ)) ,
which is an upper bound for ψ since S ⊆ S′.
We now prove an auxiliary lemma, which will be used in the proof of Lemma 3.4.
Lemma 4.1. Let α, β be two non-zero vectors of Rn. Then
∣∣∣∣
α
|α| − β
|β|
∣∣∣∣ ≤ 2 |α − β|
|α| .
Proof. This is an easy application of the triangle inequality. We have
∣∣∣∣
α
|α| − β
|β|
∣∣∣∣ = 1
|α|
∣∣∣∣α − |α|
|β| β
∣∣∣∣
and ∣∣∣∣α − |α|
|β| β
∣∣∣∣ ≤ | α − β| +
∣∣∣∣β − |α|
|β| β
∣∣∣∣ = |α − β| + || β|−| α|| ≤ 2|α − β|.
Next, we prove Lemma 3.4. 8Proof of Lemma 3.4. Without loss of generality we assume that |β1| = max( |βi| i = 1 , 2 . . . d ). Define ξi = βi
β1
for i = 2 , 3 . . . d , so that |ξi| ≤ 1 for all i. By Dirichlet’s theorem (Theorem 3.3) there exist 1 ≤ q ≤ Hd−1 and pi ∈ Z for i = 2 , 3, . . . , d such that
∣∣∣∣ξi − pi
q
∣∣∣∣ ≤ 1
qH for i = 2 , 3, . . . , d. (4.4) Let a = ( q, p 2, . . . , p d). We have
∣∣∣pi
q
∣∣∣ ≤ 1 + 1
qH
≤ 2, so that |pi| ≤ 2q. Hence
|a|2 = q2 + p22 + . . . + p2
d
≤ (4 d − 3) q2 =⇒ | a| ≤ (4 d − 3) 1/2q ≪ Hd−1. (4.5) Now let d = β1
q
a, thus d
|d|
= a
|a|
, so that by Lemma 4.1 and (4.4)
∣∣∣∣β − a
|a|
∣∣∣∣ =
∣∣∣∣β − d
|d|
∣∣∣∣ ≤ 2|β − d|= 2 |β1|
( d∑
i=2
(
ξi − pi
q
)2)1/2
≤ 2√d − 1
qH .
Since |a| ≤ (4 d − 3) 1/2q, it follows that
∣∣∣∣β − a
|a|
∣∣∣∣ ≤ 2(4 d2 − 7d + 3) 1/2
|a|H ≪ 1
|a|H . (4.6) Both implied constants in (4.5) and (4.6) are absolute.
5 Rational quotients: proof of Theorem 1.5
We now prove a generalisation of Lemma 3.4 for the case when the direction of the segment has rational quotients (recall Definition 1.4). We focus on the case when the number of rational quotients s is less than or equal to d − 2. For s = d − 1 there exists a non-zero k such that kβ ∈ Zd.Then by Lemma 3.1
ψ ≪β κd(R)(1 + h).
Arguing as in the proof of Lemma 3.2 we find that
h = 2 R(sin( θ1/4) 2 − sin( θ2/4) 2) = 2 R(sin( θ2/4) θ
4 + O(θ2)) ∼ Rθ as θ → 0Thus, this gives the same bound as Theorem 1.5, and no Diophantine approximation is needed.
Lemma 5.1. Let β ∈ Rd, |β| = 1 , have s rational quotients, 1 ≤ s ≤ d − 2, and let H ≥ 1 be an integer. Then there exists a ∈ Zd such that
|a| ≪ β Hd−1−s and
∣∣∣∣
a
|a| − β
∣∣∣∣ ≪β
1
|a|H ,
where ≪β is understood with respect to H.
9Proof. Let k ∈ R \ { 0} such that kβ has s + 1 rational coordinates. Define ξi = kβ i for i = 1 , 2, . . . d ,so that we have |ξi| ≤ | k| for all i. Without loss of generality assume that ξi are rational for
i = 1 , 2, . . . , s + 1. Then
ξi = mi
ni
∈ Q for 1 ≤ i ≤ s + 1 .
By Dirichlet’s theorem (Theorem 3.3) there exist 1 ≤ q′ ≤ Hd−1−s and pi ∈ Z such that
∣∣∣∣ξi − pi
q′
∣∣∣∣ ≤ 1
q′H for s + 2 ≤ i ≤ d. (5.1) Let m = ∏s+1
i=1
ni, q = q′m, and a =
( qm 1
n1
, qm 2
n2
. . . , qm s+1
ns+1
, mp s+2 , . . . , mp d
)
∈ Zd. Then
|a| = q
( s+1 ∑
i=1
ξ2
i
+
d
∑
i=s+2
p2
i
q′2
)1/2
≤ (1 + |k|)dq ≤ (1 + |k|)dmH d−1−s. (5.2) Now let d = 1
kq
a, then d
|d|
= a
|a|
, so that by Lemma 4.1 and (5.1)
∣∣∣∣β − a
|a|
∣∣∣∣ =
∣∣∣∣β − d
|d|
∣∣∣∣ ≤ 2|β − d|
= 2
|k|
( s+1 ∑
i=2
(
ξi − mi
ni
)2
+
d
∑
i=s+2
(
ξi − pi
q′
)2)1/2
≤ 2
|k|√d − 1 − s
q′
1
H .
Therefore, using (5.2), we obtain
∣∣∣∣β − a
|a|
∣∣∣∣ ≤ 2(1 + |k|)dm √(d − 1 − s)
|k|| a|H ≪β
1
|a|H .
This enables us to prove Theorem 1.5, in essentially the same way as Theorem 1.3.
Proof of Theorem 1.5. In view of the paragraph at the beginning of the section we can assume that 1 ≤ s ≤ d − 2. Given a segment with direction β having s rational quotients, we choose a ∈ Zd
satisfying Lemma 5.1, for some H to be determined. Then
∣∣∣∣β − a
|a|
∣∣∣∣ = 2 sin( φ/ 2) ∼ φ as φ → 0.
Later we will choose H so that θ → 0 implies φ → 0, in which case Lemma 3.2 gives
ψ ≪ κd(R)
(
1 + R|a|θ + R|a|
∣∣∣∣β − a
|a|
∣∣∣∣)
as θ → 0.
From Lemma 5.1, it follows that
|a|θ + |a|
∣∣∣∣β − a
|a|
∣∣∣∣ ≪β θH d−1−s + 1
H . (5.3) 10 Setting H =
⌈
θ −1
d−s
⌉
in (5.3), so that H = O
(
θ −1
d−s
)
, we obtain
|a|θ + |a|
∣∣∣∣β − a
|a|
∣∣∣∣ ≪β θ 1
d−s
.
Therefore,
ψ ≪β κd(R)
(
1 + Rθ 1
d−s
)
as θ → 0.
Finally, for our choice of H we obtain 2 sin( φ/ 2) =
∣∣∣∣β − a
|a|
∣∣∣∣ ≪β
1
|a|H ≤ 1
H ∼ θ 1
d−s
→ 0 as θ → 0.
Thus φ → 0 as θ → 0 since 0 ≤ φ ≤ π.
Acknowledgements
Thanks to Riccardo Maffucci for being my supervisor in this project, who has been indis-pensable in every sense. Thanks also to the LMS and Oxford University Math Institute for providing funding for this project. I would also like to thank Ze´ ev Rudnick and R. Heath-Brown for generously answering my questions about their work.
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13799 | https://tutoring.k12.com/articles/subtracting-vectors-simplified-beginners-guide/ | Subtracting Vectors: Simplified Guide for Beginners
Vector subtraction is a crucial concept in linear algebra, helping us understand how to combine and analyze quantities that have both direction and magnitude. If you’re just stepping into the world of vectors, worry not—this guide is here to simplify the topic for you. Whether you’re exploring physics, engineering, or math itself, subtracting vectors is a skill you’ll encounter often. Let’s break it down step-by-step and show you how manageable it can be!
What Are Vectors?
Before we talk about subtracting vectors, let’s first define what a vector is. A vector is a mathematical entity that represents both a magnitude (size) and a direction. Think of it as an arrow—it has a length that tells you how far to go and a direction that shows you where to aim.
Examples of Vectors in Everyday Life:
What Is Vector Subtraction?
Vector subtraction is the process of finding the difference between two vectors. It’s like asking, “If we remove the influence of one vector from another, what’s left?” This operation helps us understand relative motion, opposing forces, and more.
How to Subtract Vectors?
To subtract one vector from another, you reverse the direction of the vector you’re subtracting and then add it to the first vector. Why reverse it? Because subtraction can be viewed as the addition of a negative vector.
Let’s explore this concept visually and numerically!
Notation:
If we have two vectors A and B, the subtraction is written as:
A – B
Steps to Subtract Vectors:
You can also work it out with math, focusing on components (i.e., the x and y parts of the vector).
For example, if:
Vector A = (3, 4)
Vector B = (1, 2)
Subtract their components:
A – B = (3 – 1, 4 – 2) = (2, 2)
Real-Life Examples of Subtracting Vectors
Here are some simple, relatable scenarios to make the concept clearer:
1. Walking the Dog
Imagine you walk your dog 5 steps forward, but your dog stubbornly pulls you 2 steps backward
Subtracting vectors shows your effective movement:
A – B = 5 – 2 = 3 steps forward
2. Tug-of-War
Two teams are pulling a rope in opposite directions during a tug-of-war.
To find the net pull (resulting force):
A – B = 100 – (-70) = 170 Newtons to the right
3. Airplane and Wind
A plane is flying eastward at 500 mph, but a 50 mph wind is blowing westward.
Resultant velocity:
A – B = 500 – 50 = 450 mph East
4. Rowing Against a Current
You’re rowing your boat at 6 mph, but a 2 mph river current flows against you.
Subtract the vectors:
A – B = 6 – 2 = 4 mph forward
These examples show how vector subtraction can be applied to problems in physics, movement, and forces.
Benefits of Learning Vector Subtraction
Understanding how to subtract vectors is essential for tackling a wide range of real-world problems. Some of the benefits include:
Whether you’re an aspiring engineer, scientist, or just a math enthusiast, mastering vector operations is a step forward toward success.
Take Your Learning Further
Still want to practice and refine your skills? Check out this interactive resource from Khan Academy that includes exercises and visual examples of vector subtraction. It’s a great way to deepen your understanding and practice hands-on!
And if you’re looking for 1-on-1 help, the expert math tutors at K12 Tutoring are here to guide you every step of the way! Whether you’re dealing with vectors or any other tricky topic, our tutors make math approachable and fun.
Wrapping Up
While subtracting vectors may seem like an intimidating topic at first, breaking it down into manageable pieces makes it far friendlier. By visualizing the process and applying practical examples, you’ll gain a deeper understanding of this essential concept in linear algebra.
And if you need extra support, don’t forget—K12 Tutoring’s experts are ready to help you conquer math challenges.
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