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y satisfy the s 2 = 1, sx = −xs, sˆ y = −ˆ yx − xˆ ys, ˆ y = 1 − 2cs. The algebra generated by s, x, yˆ with these relations is the rational Cherednik algebra H1,c(Z2, C) with the action of Z2 on C is given by z → −z. 7.12. Affine and extended affine Weyl groups. Let R = be a root system with respect to a nondegenerate symmetric bilinear form (·, ) on Rn . We will assume that R is reduced. Let {αi}n i=1 ⊂ R be the set of simple roots and R+ (respectively R−) be the set of positive (respectively negative) roots. The coroots are denoted by α∨ = 2α/(α, α). � Let Q∨ = ∨ the coweight lattice, where ∨, αj ) = δij . Let θ be the maximal positive root, ωi and assume that the bilinear form is normalized by the condition (θ, θ) = 2. Let W be the Weyl group which is generated by the reflections sα (α ∈ R). ∨’s are the fundamental coweights, i.e., (ωi i be the coroot lattice and P ∨ = Zα∨ Zωi n i=1 n i=1 � · {α} ⊂ Rn By definition, the affine root system is Ra = {α˜ = [α, j] ∈ Rn × R| where α ∈ R, j ∈ Z}. The set of positive affine
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, j] ∈ Rn × R| where α ∈ R, j ∈ Z}. The set of positive affine roots is Ra = {[α, j] [−θ, 1]. We will identify α ∈ R with ˜α = [α, 0] ∈ Ra . + | j ∈ Z>0} ∪ {[α, 0] | α ∈ R+}. Define α0 = For an arbitrary affine root ˜α = [α, j] and a vector ˜z = [z, ζ] ∈ Rn × R, the corresponding affine reflection is defined as follows: sα˜(˜z) = z˜ − 2 (z, α) (α, α) α˜ = z˜ − (z, α∨) ˜α. The affine Weyl group W a is generated by the affine reflections {sα˜ | α˜ ∈ R�+}, and we have an isomorphism: W a ∼ = W � Q∨, where the translation α∨ ∈ Q∨ is naturally identified with the composition s[−α,1]sα ∈ W a. Define the extended affine Weyl group to be W a = W � P ∨ acting on Rn+1 via b(˜z) = ext [z, ζ − (b, z)] for ˜z = [z, ζ], b ∈ P ∨. Then W a ⊂ W . Moreover, W a is a normal subgroup of a ext and W /W a = P ∨/Q∨. The latter group can be identified with the group Π = {πr} of W a a permuting simple affine roots under their action in Rn+1 . It is a normal the
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W a a permuting simple affine roots under their action in Rn+1 . It is a normal the elements of W commutative subgroup of Aut = Aut(Dyna) (Dyna denotes the affine Dynkin diagram). The quotient Aut/Π is isomorphic to the group of the automorphisms preserving α0, i.e. the group AutDyn of automorphisms of the finite Dynkin diagram. ext a ext ext 7.13. Cherednik’s double affine Hecke algebra of a root system. In this subsection, we will give an explicit presentation of Cherednik’s DAHA for a root system, defined in Example 7.18. This is done by giving an explicit presentation of the corresponding braid group (which is called the elliptic braid group), and then imposing quadratic relations on the generators corresponding to reflections. For a root system R, let m = 2 if R is of type D2k, m = 1 if R is of type B2k, Ck, and otherwise m = Π . Let mij be the number of edges between vertex i and vertex j in the 58 | | affine Dynkin diagram of Ra . Let Xi (i = 1, . . . , n) be a family of pairwise commutative and algebraically independent elements. Set n � X[b,j] = Xi �i qj, where b = n � �iωi ∈ P, j ∈ Z/mZ. i=1 i=1 ext For an element wˆ ∈ W a , we can define an action on these X[b,j] by ˆwX[b,j] = Xwˆ[b,j
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can define an action on these X[b,j] by ˆwX[b,j] = Xwˆ[b,j]. Definition 7.24 (Cherednik). The double affine Hecke algebra (DAHA) of the root system R, denoted by HH, is an algebra defined over the field Cq,t = C(q1/m, t1/2), generated by Ti, i = 0, . . . , n, Π, Xb, b ∈ P , subject to the following relations: i (1) TiTj Ti · · · = Tj TiTj · · · , mij factors each side; (2) (Ti − ti)(Ti + t−1) = 0 for i = 0, . . . , n; (3) πTiπ−1 = Tπ(i), for π ∈ Π and i = 0, . . . , n; (4) πXbπ−1 = Xπ(b), for π ∈ Π, b ∈ P ; 1 , if i > 0 and (b, α∨ − (5) TiXbTi = XbXα i (6) T0XbT0 = Xb−α0 if (b, θ) = −1; T0Xb = XbT0 if (b, θ) = 0. Here ti are parameters attached to simple affine roots (so that roots of the same length i ) = 1; TiXb = XbTi, if i > 0 and (b, αi ∨) = 0; give rise to the same parameters). The degenerate double affine Hecke algebra (trigonometric Cherednik algebra) HHtrig is i=1(b, α∨)yi + u ext generated by the group algebra of W a , Π and pairwise commutative y˜b =
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�)yi + u ext generated by the group algebra of W a , Π and pairwise commutative y˜b = for ˜b = [b, u] ∈ P × Z, with the following relations: � n siyb − ysi(b)si = −ki(b, αi ∨), for i = 1, . . . , n, s0yb − ys0(b)s0 = k0(b, θ), πrybπr −1 = yπr (b) for πr ∈ Π. Remark 7.25. This degeneration can be obtained from the DAHA similarly to the case of A1, which is described above. 7.14. Algebraic flatness of Hecke algebras of polygonal Fuchsian groups. Let W be the Coxeter group of rank r corresponding to a Coxeter datum: mij (i, j = 1, . . . , r, i =� j), such that 2 ≤ mij ≤ ∞ and mij = mji. So the group W has generators si i = 1, . . . , r, and defining relations s 2 i = 1, (sisj )mij = 1 if mij =� ∞. It has a sign character ξ : W → {±1} given by ξ(si) = −1. Denote by W+ the kernel of ξ (the even subgroup of W ). It is generated by aij = sisj with relations: = a−1 ji , aij ajkaki = 1, We can deform the group algebra C[W ] as follows. Define the algebra A(W ) with invertible generators si, and tij,k, i, j = 1, . . . , r, k ∈ Zmij for (i, j) such that mij < ∞ and defining relations mij aij =
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i, j) such that mij < ∞ and defining relations mij aij = 1. aij tij,k = t−1 ji,−k, 2 = 1, si mij � [tij,k, ti�j�,k� ] = 0, sptij,k = tji,ksp, (sisj − tij,k) = 0 if mij < ∞. k=1 Notice that if we set tij,k = exp(2πki/mij ), we get C[W ]. 59 Define also the algebra A+(W ) over R := C[tij,k] (tij,k = t−1 ) by generators aij , i = j ji,−k (aij = a−1 ji ), and relations mij � (aij − tij,k) = 0 if mij < ∞, aij ajpapi = 1. k=1 If w is a word in letters si, let Tw be the corresponding element of A(W ). Choose a function w(x) which attaches to every element x ∈ W , a reduced word w(x) representing x in W . Theorem 7.26 (Etingof, Rains, [ER]). set in A(W ) as a left R-module. (i) The elements Tw(x), x ∈ W , form a spanning (ii) The elements Tw(x), x ∈ W+, form a spanning set in A+(W ) as a left R-module. (iii) The elements Tw(x), x ∈ W , are linearly independent if W has no finite parabolic subgroups of rank 3. Proof. We only give the proof of (i). Statement (ii) follows from (i). Proof of (iii), which is quite nontrivial, can be found in [ER] (it uses the geometry of constructible sheaves on the
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Proof of (iii), which is quite nontrivial, can be found in [ER] (it uses the geometry of constructible sheaves on the Coxeter complex of W ). Let us write the relation as a deformed braid relation: mij � (sisj − tij,k) = 0 k=1 sj sisj . . . + S.L.T. = tij sisj si . . . + S.L.T., where tij = (−1)mij +1tij,1 · · · tij,mij , S.L.T. mean “smaller length terms”, and the products on both sides have length mij . This can be done by multiplying the relation by sisj · · · (mij factors). Now let us show that Tw(x) span A(W ) over R. Clearly, Tw for all words w span A(W ). So we just need to take any word w and express Tw via Tw(x). It is well known from the theory of Coxeter groups (see e.g. [B]) that using the braid relations, one can turn any non-reduced word into a word that is not square free, and any reduced expression of a given element of W into any other reduced expression of the same element. Thus, if w is non-reduced, then by using the deformed braid relations we can reduce Tw to a linear combination of Tu with words u of smaller length than w. On the other hand, if w is a reduced expression for some element x ∈ W , then using the deformed braid relations we can reduce Tw to a linear combination of Tu with u shorter than w, and Tw(x). Thus Tw(x) � are a spanning set. This proves (i). Thus, A+(W ) is a “deformation” of C[W+] over R, and similarly A(W ) is a
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is a “deformation” of C[W+] over R, and similarly A(W ) is a “twisted deformation” of C[W ]. Now let Γ = Γ(m1, . . . , mr), r ≥ 3, be the Fuchsian group defined by generators cj , j = 1, . . . , r, with defining relations c m j j = 1, r � cj = 1. Here 2 ≤ mj < ∞. j=1 60 � Suppose Γ acts on H where H is a simply connected complex Riemann surface as in Section 7.7. We have the Hecke algebra of Γ, Hτ (H, Γ), defined by the same (invertible) generators cj and relations � � r (cj − exp(2πik/nj )qjk) = 0, cj = 1, k j=1 where qjk = exp(τjk). We saw above (Theorem 7.15) that if τjk’s are formal, the algebra Hτ (Γ, H) is flat in τ if |Γ| is infinite (i.e., H is Euclidean or hyperbolic). Here is a much stronger non-formal version of this theorem. Theorem 7.27. The algebra Hτ (Γ, H) is free as a left module over R := C[q±1] if and only if (1 − 1/mj ) ≥ 2 (i.e., H is Euclidean or hyperbolic). j � jk Proof. Let us consider the Coxeter datum: mij , i, j = 1, . . . , r, such that mi,i+1 := mi (i ∈ Z/rZ), and mij = ∞ otherwise. Suppose the corresponding Coxeter group is
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mi,i+1 := mi (i ∈ Z/rZ), and mij = ∞ otherwise. Suppose the corresponding Coxeter group is W . Then we can see that Γ = W+. Notice that the algebra Hτ (Γ, H) for genus 0 orbifolds is the algebra A+(W ), i.e., we have Hτ (Γ, H) = A+(W ). The condition (1 − 1/mj ) ≥ 2 is equivalent to the condition that W has no finite j parabolic subgroups of rank 3. From Theorem 7.26 (ii) and Theorem 7.15, we can see that � A+(W ) is free as a left module over R. We are done. � 7.15. Notes. Section 7.8 follows Section 6 of the paper [EOR]; Cherednik’s definition of the double affine Hecke algebra of a root system is from Cherednik’s book [Ch]; Sections 7.7 and 7.14 follow the paper [ER]; The other parts of this section follow the paper [E1]. 61 MIT OpenCourseWare http://ocw.mit.edu 18.735 Double Affine Hecke Algebras in Representation Theory, Combinatorics, Geometry, and Mathematical Physics Fall 2009 For information about citing these materials or our Terms of Use, visit: http://ocw.mit.edu/terms.
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3.044 MATERIALS PROCESSING LECTURE 4 General Heat Conduction Solutions: ∂T = ∇ · k∇T, T (¯x, t) ∂t Trick one: steady state ∇2T = 0, T (x) Trick two: low Biot number ∂T = α h(Ts − Tf ), T (t) ∂t Low Biot Number Solutions: Newtonian Heating / Cooling Global Heat Balance: qconv = qlost A h(T − Tf ) = −ρ cp ∂T T − Tf = ∂T ∂t −hA ρcpV V ∂t ln(T − Tf ) = −hA ρ cp V @t = 0, T = Ts t + C Date: February 21st, 2012. ln(Ts − Tf ) = C 1 2 LECTURE 4 ln T − Tf Ts − Tf T − Tf Ts − Tf = −hA ρ cp V t −hA ρ cp V = e t Transient Heat Conduction: depends on position and time ∂T ∂t = α ∇2 T You should know: 1) Some common solutions for simple geometries 2) Where to find solutions 3) How to build up complex solutions using simple solutions Semi-Infinite Solid - constant T1 at surface - initially T0 everywhere T − T0 T1 − T0 x = erfc √ 2 αt 3.044 MATERIALS PROCESSING 3 erf(z) = (cid:90) z 0 e−x2dx erfc = 1 − erf T (x) = (T1 − T0) erfc (cid:18) (cid:19) x √ α t 2 + T0 T − T0 T1 − T0 T − T0 T1 − T0 (−1) = erfc (cid:18) = er
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1 − T0 T − T0 T1 − T0 (−1) = erfc (cid:18) = erfc (cid:18) (cid:19) x √ α t x (cid:19) √ α t 2 2 (−1) 4 LECTURE 4 − (T − T0) T0 − T T1 − T0 T − T1 T1 − T0 T − T1 T0 − T1 (cid:19) (cid:18) x = erf √ 2 α t x (cid:18) = −erf √ 2 α t x (cid:19) (cid:18) = erf √ 2 α t (cid:19) (cid:18) hx k (cid:19) + h2k α t Semi-Infinite Solid - convection at surface: qlost = h(T − Tf ) Θ = ERF C (cid:19) (cid:18) x √ 2 α t − EXP Where to find these solutions: - Carslaw & Jaeger - Crank Dimensionless Numbers: x = erfc √ 2 αt T − T0 T1 − T0 T − T0 = Θ T1 − T0 x χ = L x = Lχ α t τ = L2 L2T t = α ⎛ Θ = erfc L ⎝ Lχ α L2 τ α 2 (cid:18) χ (cid:19) Θ = erfc √ 2 τ · ERF C (cid:18) x √ + 2 α t h √ k (cid:19) α t ⎞ ⎠ 3.044 MATERIALS PROCESSING 5 MIT OpenCourseWare http://ocw.mit.edu 3.044 Materials Processing Spring 2013 For information about citing these materials or our Terms of Use, visit
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Finite Element Modeling of the Detachment of Soft Adhesives Stick-slip phenomena and Schallamach waves captured using reversible cohesive elements 1 Evelyne Ringoot • • • • BSc in Engineering Sciences at VUB Brussels 2018 Msc in Civil Engineering at VUB Brussels 2020 o Specialization in geomechanics and numerical methods Visiting Student at École polytechnique fédérale de Lausanne Visiting Student at Massachusetts Institute of Technology 2 Soft Adhesives And the remarkable reversible capacities of natural adhesives How to explain reattachment and reversible adhesion? © National Geographic Society. All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. Stefan Sirucek (2014) How gecko’s turn their stickiness on and off in ‘National Geographic’. 3 Medical: tissue repair, wound scaffolds or drug patches High-precision non-damaging soft grippers © UC San Diego Jacobs School of Engineering. All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. Soft Adhesive Applications © Karp Laboratory. All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. © The European Space Agency. All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. Climbing robots for dangerous environments UC San Diego Jacobs School of Engineering (2018), Tolley Gecko Gripper on Flickr, consulted in Sept 2020 on https://www.flickr.com/photos/jsoe/albums/72157695462669655/with/40449351705/ The European Space Agency (2014), Wall-crawling gecko robots can stick in space too, consulted in Sept 2020 http://www.esa.int/Enabling_Support/Space_Engineering_Technology/Wall-crawling_gecko_robots_can_stick_in_space_too The Karplab (2014), Worm-Inspired Glue Mends Broken Hearts, consulted on Sept 2020 on https://www.karplab.net
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plab (2014), Worm-Inspired Glue Mends Broken Hearts, consulted on Sept 2020 on https://www.karplab.net/portfolio-item/worm-inspired-glue-mends-broken-hearts 4 Experimental observations Research questions in mechanics of solids: how to explain, predict and influence physical realities? Analytical theory Cohesive elements represent surface strength assumptions Numerical solutions 5 Finite Element Models of solid deformation Differential equations governing the conservation of mass and momentum: + constitutive equations linking stress induced by forces to strain encountered by the material +Boundary conditions on the stress or strain state applied on the borders of the material 6 Commercial Finite Element Models software ANSYS ABAQUS NX NASTRAN Or code developed in research groups: Akantu 7 The detachment and re-attachment of adhesive with multiple layers when loaded parallel to their substrate 8 Adapting a FEM framework allowed to numerically replicate a physical phenomena that is still not fully understood: Soft Adhesive detachment Experiment Simulation Attached Detached Attached Detached 9 MIT OpenCourseWare https://ocw.mit.edu 18.085 Computational Science and Engineering I Summer 2020 For information about citing these materials or our Terms of Use, visit: https://ocw.mit.edu/terms. 10
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MORE MATLAB INSTRUCTIONS 1. Plotting functions In this section, the basics of plotting functions in MATLAB are described. Throughout we work with the example of two functions, f (t) = (2t + 1)e−t sin(t), and g(t) = (t − 1)e−t cos(t). Step 1, Specify the domain: Functions are defined on an interval called the domain. To plot the function in MATLAB, you need to specify the domain. Every domain has a left endpoint, a, and a right endpoint, b. Of course MATLAB does not plot the value of the function at every point between the a and b, only finitely many points with a regular spacing h. The syntax to specify the domain is, >> t = a : h : b For example, to plot our function on the interval [−1, 1] with step size 0.05, the syntax is, >> t = ­1:0.05:1 just the ordered list of the One word about this. Technically x is a data type called an array: numbers a, a + h, a + 2h, . . . . The syntax for arithmetic with an array in MATLAB is different than the syntax for arithmetic with a number. Step 2, Specify the function: Here is a list of common operations used to define functions, and the corresponding syntax in MATLAB. In the list, y(t)
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define functions, and the corresponding syntax in MATLAB. In the list, y(t) and z(t) are names for functions or pieces of functions that are already specified. Operation y(t) + z(t) y(t)z(t) y(t)n y(t)/z(t) sin(y(t)) cos(y(t)) ey(t) ln(y(t)) log10(y(t)) MATLAB Syntax y + z y.* z y.^n y./z sin(y) cos(y) exp(y) log(y) log10(y) For example, if the range t has already been defined, the function (2t + 1)e−t sin(t) is specified by, >> y = ( 2 .* t + 1 ).* exp( ­1.* t ).* sin( t ) Similarly, the function (t − 1)e−t cos(t) is specified by, >> z = ( t ­ 1 ).* exp( ­1.* t ).* cos( t ) Step 3, Plot the function: The syntax to produce a 2D­plot whose domain is t and whose function is y is, >> h = plot(t,y) Note, you do not need to say “h =”, but this can be useful if you want to manipulate the plot later. MATLAB will produce the plot in a new window. Step 4, Plotting a parametrized curve; Several plots at once: MATLAB can plot a parametrized figure. For instance, for the parametrized curve (y
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plot a parametrized figure. For instance, for the parametrized curve (y, z) where y(t) = (2t + 1)e−t sin(t), z(t) = (t − 1)e−t cos(t), the syntax is, 1 >> i = plot(y,z) where y and z are specified as above. Note that when plotting parametrized curves, it is still necessary to specify the t­domain. But t doesn’t explicitly appear in the syntax of the plot. Also, MATLAB can plot several graphs (or parametrized curves) simultaneously. For simplicity, think of a graph as a parametrized curve (t, y(t)). For a number of parametrized curves, say (y1(t), z1(t)), (y2(t), z2(t)), the syntax to plot both of these curves in a single figure is, >> j = plot(y1,z1,y2,z2) Any number of curves can be plotted in a single figure: just write plot(y1, z1, . . . , yn, zn), where the functions y1, . . . , yn and z1, . . . , zn have already been specified. To simultaneously graph the functions y and z above over the interval t, the syntax is, >> k = plot(t,y,t,z) Step 5, Print or export your plot: To either print your plot or to export it as a JPEG file, click on the “File” button of the new window and
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a JPEG file, click on the “File” button of the new window and then click on “Print” or “Export” in the pop­up menu. There are other extras that you can find out by experimenting (such as adding labels to your axes). 2
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MIT OpenCourseWare http://ocw.mit.edu 18.306 Advanced Partial Differential Equations with Applications Fall 2009 For information about citing these materials or our Terms of Use, visit: http://ocw.mit.edu/terms. Discrete to Continuum Modeling. Rodolfo R. Rosales . (cid:3) MIT, March, 2001. Abstract These notes give a few examples illustrating how continuum models can be derived from special limits of discrete models. Only the simplest cases are considered, illustrating some of the most basic ideas. These techniques are useful because continuum models are often much easier to deal with than discrete models with very many variables, both conceptually and computationally. Contents 1 Introduction. 2 2 Wave Equations from Mass-Spring Systems. 2 Longitudinal Motion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 2 Nonlinear Elastic Wave Equation (for a Rod) . . . . . . . . . . . . . . . . . . . . . . 4 Example: Uniform Case . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 4 | Sound Speed . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 4 Example: Small Disturbances . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 5 | Linear Wave Equation, and Solutions . . . . . . . . . . . . . . . . . . . . . . . . 5 Fast Vibrations . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 5 | Dispersion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 6 | Long Wave Limit . . . . . . . . . . . . . . . . . . .
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. . . . . . 6 | Long Wave Limit . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 6 Transversal Motion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 6 Stability of the Equilibrium Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . 7 Nonlinear Elastic Wave Equation (for a String) . . . . . . . . . . . . . . . . . . . . . 7 Example: Uniform String with Small Disturbances . . . . . . . . . . . . . . . . . . . 7 | Uniform String Nonlinear Wave Equation. . . . . . . . . . . . . . . . . . . . . . . 7 | Linear Wave Equation. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 7 | Stability and Laplace’s Equation. . . . . . . . . . . . . . . . . . . . . . . . . . . . 8 | Ill-posed Time Evolution. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 8 General Motion: Strings and Rods . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 8 3 Torsion Coupled Pendulums: Sine-Gordon Equation. 8 Hooke’s Law for Torsional Forces . . . . . . . . . . . . . . . . . . . . . . . . . . . . 9 Equations for N torsion coupled equal pendulums . . . . . . . . . . . . . . . . . . . 10 Continuum Limit . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 10 | S
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. . . . . . . . . . . . . . . . . . . . . . . . . . . . 10 | Sine-Gordon Equation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 11 | Boundary Conditions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 11 Kinks and Breathers for the Sine Gordon Equation . . . . . . . . . . . . . . . . . . . . 11 Example: Kink and Anti-Kink Solutions . . . . . . . . . . . . . . . . . . . . . . . . 12 Example: Breather Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 13 Pseudo-spectral Numerical Method for the Sine-Gordon Equation . . . . . . . . . . . . 14 4 Suggested problems. 14 (cid:3) MIT, Department of Mathematics, room 2-337, Cambridge, MA 02139. 1 Discrete to Continuum Modeling. 2 MIT, March, 2001 | Rosales. 1 Introduction. Continuum approximations are useful in describing discrete systems with a large number of degrees of freedom. In general, a continuum approximation will not describe all possible solutions of the discrete system, but some special class that will depend on the approximations and assumptions made in deriving the continuum model. Whether or not the approximation is useful in describing a particular situation, will depend on the appropriate approximations being made. The most successful models arise in situations where most solutions of the discrete model evolve rapidly in time towards con(cid:12)gurations where the assumptions behind the continuum model apply. The basic step in obtaining a continuum model from a discrete system, is to identify some basic con(cid:12)guration (solution of the discrete model) that can be described by a few parameters. Then one assumes that the full solution of the system can be described, near every point in space and at every time, by this con(cid:12)guration | for some value of the parameters. The parameters are
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, near every point in space and at every time, by this con(cid:12)guration | for some value of the parameters. The parameters are then assumed to vary in space and time, but on scales (macro-scales) that are much larger than the ones associated with the basic con(cid:12)guration (micro-scales). Then one attempts to derive equations describing the evolution of these parameters in the macro-scales, thus averaging out of the problem the micro- scales. There is a close connection between this approach, and the"quasi-equilibrium" approximations that are often invoked to "close" continuum sets of equations derived using conservation laws. For example, when deriving the equations for Gas Dynamics in Statistical Mechanics, it is assumed that the local particle interactions rapidly exchange energy and momentum between the molecules | so that the local probability distributions for velocities take a standard form (equivalent to local thermodynamic equilibrium). What exactly makes these assumptions work (in terms of properties of the governing, micro-scale, equations) is rather poorly understood. But that they work rather well cannot be denied. In these notes we will consider examples that are rather simpler than these ones, however, where the "local con(cid:12)gurations" tend to be rather trivial. 2 Wave Equations from Mass-Spring Systems. Consider an array of bodies/particles, connected by springs, and restricted to move on a straight Longitudinal Motion. 1 line. Let the positions of the bodies be given by x = x (t), with n = 0; 1; 2; : : :, and let M n n n be the mass of the n particle. Furthermore, let the force law for the spring between particles th (cid:6) (cid:6) n and n + 1 be given by: force = f ((cid:1)x), where (cid:1)x is the distance between the particles, and 1 1 f is positive when the spring is under tension. + n 2 + n 2 2 If there are no other forces involved (e.g. no friction), the governing equations for the system are: 2 d M x = f (x x ) f (x x ) ; (2.1) +1 1 n n + n n n n
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) f (x x ) ; (2.1) +1 1 n n + n n n n(cid:0) 2 n n(cid:0) 2 2 1 1 dt (cid:0) (cid:0) (cid:0) for n = 0; 1; 2; : : : The simplest solution for this system of equations is equilibrium. In this case all the accelerations vanish, so that the particle positions are given by the series of algebraic (cid:6) (cid:6) By some device: say the bodies are sliding inside a hollow tube. 1 2 If the spring obeys Hooke’s law, then ((cid:1) ) = (cid:1) , where 0 and 0 are the n n n n n f x k x L k > L > + + + + + 2 2 2 2 2 1 1 1 1 1 (cid:0) spring constant and equilibrium length, respectively. (cid:16) (cid:17) Discrete to Continuum Modeling. 3 MIT, March, 2001 | Rosales. equations 0 = f (x x ) f (x x ) : (2.2) +1 1 + n n n n(cid:0) n n(cid:0) 2 2 1 1 (cid:0) (cid:0) (cid:0) This is the basic con(cid:12)guration (solution) that we will use in obtaining a continuum approximation. Note that this is a one parameter family: if the forces are monotone functions of the displacements (cid:1)x, then once any one of them is given, the others follow from (2.2). Before proceeding any further, it is a good idea to non-dimensionalize the equations. We will assume that: A. All the springs are roughly similar, so that we can talk of a typical spring force f , and a typical spring length L. Thus we can write 1 1 (cid:1)x f ((cid:1)x) = f F ; (2.3) + + n n 2 2 L (cid:18) (cid:19) where F is a
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; (2.3) + + n n 2 2 L (cid:18) (cid:19) where F is a non-dimensional mathematical function, of O(1) size, and with O(1) deriva- + n 2 1 tives. A further assumption is that F changes slowly with n, so that two nearby springs + n 2 1 are nearly equal. Mathematically, this is speci(cid:12)ed by stating that: 1 F ((cid:17) ) = F ((cid:15)(n + 1=2); (cid:17) ) ; (2.4) + n 2 where 0 < (cid:15) 1, and F is a "nice" (mathematical) function of its two variables. (cid:28) B. All the particles have roughly the same mass m, and their masses change slowly with n, so that we can write: M = m M ((cid:15) n) ; (2.5) n where M is a nice mathematical function, with O(1) size, and with O(1) derivatives. Remark 2.1 Why do we need these assumptions? This has to do with the questions of validity, discussed in the introduction. Suppose that these hypothesis are violated, with the masses and springs jumping wild ly in characteristics. Then the basic con(cid:12)guration described by (2.2) wil l stil l be a solution. However, as soon as there is any signi(cid:12)cant motion, neighboring parts of the chain wil l respond very di(cid:11)erently, and the solution wil l move away from the local equilibrium implied by (2.2). There is no known method to, generical ly, deal with these sort of problems | which turn out to be very important: see remark 2.2. From the assumptions in A and B above, we see that: Changes in the mass-spring system occur over length scales ‘ = L=(cid:15): (2.6) Using this scale to non-dimensionalize space, namely: x = ‘ X ; and a yet to be speci(cid:12)ed time n n scale (cid:28) to non-dimensionalize time, namely: t = (cid:28) T ;
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12)ed time n n scale (cid:28) to non-dimensionalize time, namely: t = (cid:28) T ; the equations become: 2 2 d (cid:15) f (cid:28) X X X X +1 1 n n n n(cid:0) 1 1 M ((cid:15)n) X = F F : (2.7) n + (cid:0) (cid:0) 2 2 2 n n(cid:0) dT m L (cid:15) (cid:15) (cid:0) (cid:18) (cid:18) (cid:19) (cid:18) (cid:19)(cid:19) A and B above also imply that, for the solution in (2.2), the inter-particle distance x x varies +1 n n (cid:0) slowly | an O((cid:15)) fractional amount per step in n. Thus we propose solutions for (2.7) of the form: and X = X (s; t) is some smooth function of its arguments. X (t) = X (s ; t) ; where s = n (cid:15) ; (2.8) n n n Discrete to Continuum Modeling. 4 MIT, March, 2001 | Rosales. Substituting (2.8) into (2.7), and using (2.4) and (2.5), we obtain 2 2 2 @ (cid:15) f (cid:28) @ @ 2 M (s) X = F s; X + O((cid:15) ) : (2.9) 2 @T m L @ s @ s ! ! Here we have used that: X X @ 1 X X @ 1 +1 1 n n 2 n n(cid:0) 2 (cid:0) (cid:0) = X (s + (cid:15); t) + O((cid:15) ) and = X (s (cid:15); t) + O((cid:15) ) ; (cid:15) @ s 2 (cid:15) @ s 2 (cid:
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); t) + O((cid:15) ) ; (cid:15) @ s 2 (cid:15) @ s 2 (cid:0) with a similar formula applying to the di(cid:11)erence F F . 1 1 + n n(cid:0) 2 2 (cid:0) Equation (2.9) suggests that we should take m L (cid:28) = ; (2.10) 2 s (cid:15) f for the un-speci(cid:12)ed time scale in (2.7). Then equation (2.9) leads to the continuum limit ap- proximations (valid for 0 < (cid:15) 1) (cid:28) 2 @ @ @ M (s) X = F s; X : (2.11) 2 @T @ s @ s ! The mass-spring system introduced in equation (2.1) can be thought of as a simple model for an elastic rod under (only) longitudinal forces. Then we see that (2.11) is a model (nonlinear wave) equation for the longitudinal vibrations of an elastic rod, with s a lagrangian coordinate for the points in the rod, M = M (s) the mass density along the rod, and X giving the position of the point s as a function of time, and F a function characterizing the elastic response of the rod. Of course, in practice F must be obtained from laboratory measurements. Remark 2.2 The way in which the equations for nonlinear elasticity can be derived for a crystal line solid is not too di(cid:11)erent from the derivation of the wave equation (2.11) for longitudinal vibrations. 3 Then a very important question arises (see (cid:12)rst paragraph in section 1): What important behaviors are missed due to the assumptions in the derivation? How can they be modeled? In particular, what happens if there are "defects" in the crystal structure (see remark 2.1)? These are al l very important, and open, problems of current research interest. Example 2.1 Uniform Rod. If al l the springs and al l the particles are equal, then we can take M 1 and F is independent of s. Furthermore, if we take L to be the (common) equilibrium length of
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equal, then we can take M 1 and F is independent of s. Furthermore, if we take L to be the (common) equilibrium length of the springs, we then have (cid:17) 2 2 @ @ @ @ @ 2 X = F X = c X X ; (2.12) 2 2 @T @ s @ s @ s @ s ! ! 2 2 where c = c ((cid:17) ) = dF =d(cid:17) ((cid:17) ) > 0, and F (1) = 0 (equilibrium length). The unperturbed "rod" cor- responds to X s, while X (cid:11) s corresponds to the rod under uniform tension ((cid:11) > 1), or com- pression ((cid:11) < 1). Also, note that c is a (non-dimensional) speed | the speed at which elastic (cid:17) (cid:17) disturbances along the rod propagate: i.e. the sound speed. 3 At least qualitatively, though it is technically far more challenging. Discrete to Continuum Modeling. 5 MIT, March, 2001 | Rosales. Example 2.2 Small Disturbances. Consider a uniform rod in a situation where the departures from uniform equilibrium are smal l. That is @X=@ s (cid:11), where (cid:11) is a constant. Then equation (2.12) can be approximated by the linear wave equation (cid:25) 2 X = c X ; (2.13) T T ss where c = c((cid:11)) is a constant. The general solution to this equation has the form where g and h are arbitrary functions. This solution clearly shows that c is the wave propagation X = g (s c T ) + h(s + c T ) ; (2.14) (cid:0) velocity. Remark 2.3 Fast vibrations. The vibration frequency for a typical mass m, attached to a typical spring in the chain, is: f 1 ! = = : (2.15) s m L (cid:15)(cid:28) This corresponds to a time scale much shorter than the one involved in the solution in (2.8
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m L (cid:15)(cid:28) This corresponds to a time scale much shorter than the one involved in the solution in (2.8{2.11). What role do the motions in these scales play in the behavior of the solutions of (2.1), under the assumptions made earlier in A and B? For real crystal lattices, which are de(cid:12)nitely not one dimensional (as the one in (2.1)) these fast time scales correspond to thermal energy (energy stored in the local vibrations of the atoms, relative to their equilibrium positions). It is believed that the nonlinearities in the lattice act so as to randomize these vibrations, so that the energy they contain propagates as heat (di(cid:11)uses). In one dimension, however, this does not general ly happen, with the vibrations remaining coherent enough to propagate with a strong wave component. The actual processes involved are very poorly understood, and the statements just made result, mainly, from numerical experiments with nonlinear lattices. Just to be a bit more precise: consider the situation where al l the masses are equal | M = m n for al l n, and al l the springs are equal and satisfy Hooke’s law (linear elasticity): 1 (cid:1)x f ((cid:1)x) = k((cid:1)x L) = f 1 ; (2.16) + n 2 (cid:0) (cid:0) L (cid:18) (cid:19) where k is the spring constant, L is the equilibrium length, and f = k L: Then equation (2.1) takes the form 2 d 2 x = ! (x 2x + x ) ; (2.17) +1 1 2 n n n n(cid:0) dt (cid:0) where ! is as in (2.15). Because this system is linear, we can write its general solution as a linear superposition of eigenmodes, which are solutions of the form 4 x = exp(i (cid:20) n i (cid:27) t) ; where (cid:27) = 2 ! sin and < (cid:20) < is a constant. (2.18) n (cid:20) (cid:0) (cid
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sin and < (cid:20) < is a constant. (2.18) n (cid:20) (cid:0) (cid:6) (cid:0) 1 1 2 (cid:18) (cid:19) These must be added to an equilibrium solution x = (cid:11) L n = s , where (cid:11) > 0 is a constant. n n 4 Check that these are solutions. Discrete to Continuum Modeling. 6 MIT, March, 2001 | Rosales. Relative to the mean position s along the lattice, each solution in (2.18) can be written as n (cid:20) x = exp(i s i (cid:27) t) : n n (cid:11)L (cid:0) Thus we see that it represents a wave of wavelength (cid:21) = 2(cid:25)(cid:11)L=(cid:20), and speed (cid:11)L(cid:27) 2(cid:11)L! (cid:20) 2c (cid:20) c = = sin = sin (2.19) w (cid:20) (cid:20) 2 (cid:20) 2 (cid:6) (cid:18) (cid:19) (cid:18) (cid:19) propagating along the lattice | where c = (cid:11)L! is a speed. Note that the speed of propagation is a function of the wave-length | this phenomenon is know by the name of dispersion. We also note that the maximum frequency these eigenmodes can have is (cid:27) = 2!, and corresponds to wavelengths of the order of the lattice separation. 5 In the case of equations (2.16 { 2.17) there is no intrinsic (cid:15) in the equations: it must arise from the initial conditions. That is to say: assume that the wavelength ‘ with which the lattice is excited is much larger than the lattice equilibrium separation L, i.e. ‘ L, with (cid:15) = L=‘. This corresponds to solutions (2.18) with (cid:20) smal l. In this long wave limit we see that (2.19) implies that the solutions (cid:29) have the same wave speed c = c. This
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this long wave limit we see that (2.19) implies that the solutions (cid:29) have the same wave speed c = c. This corresponds to the situation in (2.13 { 2.14). w (cid:6) It is clear that, in the linear lattice situation described above, we cannot dismiss the fast vibration excitations (with frequencies of the order of !) as constituting some sort of energy "bath" to be interpreted as heat. The energy in these vibrations propagates as waves through the media, with speeds which are of the same order of magnitude as the sound waves equation (2.13) describes. Before the advent of computers it was believed that nonlinearity would destroy the coherence of these fast vibrations. Numerical experiments, however, have shown that this is not (general ly) true for one dimensional lattices, though it seems to be true in higher dimensions. Exactly why, and 6 how, this happens is a subject of some current interest. Transversal Motion. We consider now a slightly di(cid:11)erent situation, in which the masses are allowed to move only in the direction perpendicular to the x axis. To be precise: consider a sequence of masses M in n the plane, whose x coordinates are given by x = n L. Each mass is restricted to move only in the n orthogonal coordinate direction, with y = y (t) giving its y position. The masses are connected by n n springs, with f ((cid:1)r ) the force law, where (cid:1)r = L + (y y ) is the distance between +1 + + + n n n n n 2 2 2 1 1 1 2 2 (cid:0) masses. Assuming that there are no other forces involved, the governing equations for the system q are: 2 d y y y y +1 1 n n n n(cid:0) 1 1 1 1 M y = f ((cid:1)r ) f ((cid:1)r ) ; (2.20) n n + + (cid:0) (cid:0) 2 n n n(cid:0) n
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; (2.20) n n + + (cid:0) (cid:0) 2 n n n(cid:0) n(cid:0) 2 2 2 2 1 1 dt (cid:1)r (cid:0) (cid:1)r + n 2 2 n(cid:0) for n = 0; 1; 2; : : : (you should convince yourself that this is the case). (cid:6) (cid:6) The simplest solution for this system of equations is equilibrium, with all the masses lined up horizontally y = y ; so that all the accelerations vanish. Again, one can use this (one +1 n n parameter) family of solutions to obtain a continuum approximation for the system in (2.20) | under the same assumptions earlier in A and B. The reason for the 2 relative to (2.15) is that the masses are coupled, and not attached to a single spring. 5 6 The (cid:12)rst observation of this general phenomena was reported by E. Fermi, J. Pasta and S. Ulam, in 1955: Studies of Non Linear Problems Collected Papers of Enrico Fermi. II , Los Alamos Report LA-1940 (1955), pp. 978-988 in , The University of Chicago Press, Chicago, (1965). Discrete to Continuum Modeling. 7 MIT, March, 2001 | Rosales. Remark 2.4 Stability of the Equilibrium Solutions. It should be intuitively obvious that the equilibrium solutions described above wil l be stable only if the equilibrium lengths of the springs are smaller than the horizontal separation L between the masses, 1 1 + n 2 L namely: < L. This so that none of the springs is under compression in the solution, since any + n 2 L mass in a situation where its springs are under compression wil l easily "pop" out of alignment with the others | see example 2.3. Introduce now the non-dimensional variables Y = (cid:15) y=L, X = (cid:15) x=L (note that, since x = nL, in n fact X plays here the same role that s played in the prior derivation ), and T = t=(cid
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x = nL, in n fact X plays here the same role that s played in the prior derivation ), and T = t=(cid:28) , where (cid:28) is as 7 in (2.10). Then the continuum limit for the equations in (2.20) is given by 2 @ Y @ F (X; ) @Y M (X ) = (2.21) S 2 @T @X @X ! where Y = Y (X; T ) and S 2 @Y = 1 + : v S u @X ! u t The derivation of this equation is left as an exercise to the reader. The mass-spring system introduced in (2.20) can be thought of as a simple model for an elastic string restricted to move in the transversal direction only. Then we see that (2.21) is a model (nonlinear wave) equation for the transversal vibrations of a string, where X is the longitudinal coordinate along the string position, Y is the transversal coordinate, M = M (X ) is the mass density along the string, and F = F (X; ) describes the elastic properties of the string. In 8 the non-dimensional coordinates, the (local) equilibrium length for the string is given by e = =L. S ‘ L That is, the elastic forces vanish for this length: F (X; e (X )) 0 ; where e < 1 (for stability, see remark 2.4). (2.22) ‘ ‘ (cid:17) @ We also assume that F (X; ) > 0: @ S S Example 2.3 Uniform String with Small Disturbances. Consider now a uniform string (neither M , nor F , depend on X ) in a situation where the departures from equilibrium are smal l (@Y =@X is smal l). For a uniform string we can assume M 1, and F is independent of X . Thus equation (2.21) reduces to (cid:17) 2 @ Y @ F ( ) @Y = : (2.23) S 2 @T @X @X ! S Next, for smal l disturbances we have
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: (2.23) S 2 @T @X @X ! S Next, for smal l disturbances we have 1, and (2.23) can be approximated by the linear wave equation 2 S (cid:25) 2 Y = c Y ; (2.24) T T XX where c = F (1) is a constant (see equations (2.13 { 2.14). 7 The coordinate is simply a label for the masses. Since in this case the masses do not move horizontally, can s X be used as the label. 8 S Notice that is the local stretching of the string, due to its inclination relative to the horizontal position (actual length divided by horizontal length). Discrete to Continuum Modeling. 8 MIT, March, 2001 | Rosales. Notice how the stability condition e < 1 in (2.22) guarantees that c > 0 in (2.23). If this were not ‘ 2 the case, instead of the linear wave equation, the linearized equation would have been of the form 2 Y + d Y = 0 ; (2.25) T T XX with d > 0. This is Laplace Equation, which is ill-posed as an evolution in time problem. To see this, it is enough to notice that (2.25) has the fol lowing solutions: d jkj t Y = e sin(kX ) ; for any < k < : (2.26) (cid:0) 1 1 These solutions grow arbitrarily fast in time, the fastest the shortest the wave-length ( k larger). This is just the mathematical form of the obvious physical fact that a straight string (with no bending j j strength) is not a very stable object when under compression. General Motion: Strings and Rods. If no restrictions to longitudinal (as in (2.1)) or transversal (as in (2.20)) motion are imposed on the mass-spring chain, then (in the continuum limit) general equations including both longitudinal and transversal modes of vibration for a string are obtained. Since strings have no bending strength, these equations will be well behaved only as long as the string is under tension everywhere. B
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vibration for a string are obtained. Since strings have no bending strength, these equations will be well behaved only as long as the string is under tension everywhere. Bending strength is easily incorporated into the mass-spring chain model. Basically, what we need to do is to incorporate, at the location of each mass point, a bending spring. These springs apply a torque when their ends are bent, and will exert a force when-ever the chain is not straight. The continuum limit of a model like this will be equations describing the vibrations of a rod. We will not develop these model equations here. 3 Torsion Coupled Pendulums: Sine-Gordon Equation. Consider an horizontal axle A, of total length ‘, suspended at its ends by "frictionless" bearings. Along this axle, at equally spaced intervals, there are N equal pendulums. Each pendulum consists of a rigid rod, attached perpendicularly to the axle, with a mass at the end. When at rest, all the pendulums point down the vertical. We now make the following assumptions and approximations: 1. Each pendulum has a mass . The distance from its center of mass to the axle center is L. M (cid:15) N (cid:15) (cid:15) (cid:15) 2. The axle A is free to rotate, and we can ignore any frictional forces (i.e.: they are small). In fact, the only forces that we will consider are gravity, and the torsional forces induced on the axle when the pendulums are not all aligned. 3. Any deformations to the axle and rod shapes are small enough that we can ignore them. Thus the axle and rod are assumed straight at all times. 4. The mass of the axle is small compared to M , so we ignore it (this assumption is not strictly needed, but we make it to keep matters simple). Our aim is to produce a continuum approximation for this system, as N , with everything else (cid:12)xed. ! 1 Discrete to Continuum Modeling. 9 MIT, March, 2001 | Rosales. Each one of the pendulums can be characterized by the angle (cid:18) = (cid:18) (t) that its suspending n n rod makes with the vertical direction. Each pendulum is then sub ject to three forces: (a) Gravity,
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that its suspending n n rod makes with the vertical direction. Each pendulum is then sub ject to three forces: (a) Gravity, for which only the component perpendicular to the pendulum rod is considered. 9 (b) Axle torsional force due to the twist (cid:18) (cid:18) . This couples each pendulum to the next one. +1 n n (c) Axle torsional force due to the twist (cid:18) (cid:18) . This couples each pendulum to the prior one. (cid:0) 1 n n(cid:0) (cid:0) We will assume that the amount of twist per unit length in the axle is small, so that Hooke’s law applies. Remark 3.1 Hooke’s Law for Torsional Forces. In the Hooke’s law regime, for a given (cid:12)xed bar, the torque generated is directly proportional to the angle of twist, and inversely proportional to the distance over which the twist occurs. To be speci(cid:12)c: in the problem here, imagine that a section of length (cid:1)‘ of the axle has been twisted by an amount (angle) (cid:9). Then, if T is the torque generated by this twist, one can write (cid:20) (cid:9) T = ; (3.1) (cid:1)‘ where (cid:20) is a constant that depends on the axle material and the area of its cross-section | assume that the axle is an homogeneous cylinder. The dimensions of (cid:20) are given by: mass length force area 3 [(cid:20)] = = : (3.2) (cid:2) (cid:2) 2 time angle angle (cid:2) This torque then translates onto a tangential force of magnitude F = T =L, on a mass attached to the axle at a distance L. The sign of the force is such that it opposes the twist. Let us now go back to our problem, and write the equations of motion for the N pendulums. We will assume that: The horizontal separation between pendulums is . ‘ (cid:15) N + 1 ‘ The (cid:12)rst and last pendulum are at a distance from the respective ends of the axle. (cid:15)
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1 ‘ The (cid:12)rst and last pendulum are at a distance from the respective ends of the axle. (cid:15) 2(N + 1) The tangential force (perpendicular to the pendulum rod) due to gravity on each of the masses is 1 F = M g sin (cid:18) ; where n = 1; : : : ; N : (3.3) g n (cid:0) N For any two successive masses, there is also a torque whenever (cid:18) = (cid:18) . This is generated by the +1 n n twist in the axle, of magnitude (cid:18) (cid:18) , over the segment of length ‘=(N + 1) connecting the two +1 n n (cid:0) rods. Thus each of the masses experiences a force (equal in magnitude and opposite in sign) F = (N + 1) ((cid:18) (cid:18) ) ; (3.4) +1 T n n (cid:6) (cid:0) ‘L (cid:20) where the signs are such that the forces tend to make (cid:18) = (cid:18) . Putting all this together, we obtain +1 n n the following set of equations for the angles: 1 d (cid:18) 1 (N + 1) (cid:20) 1 2 M L = M g sin (cid:18) + ((cid:18) (cid:18) ) ; (3.5) 1 2 1 2 N dt (cid:0) N ‘L (cid:0) 9 The component along the rod is balanced by the rod itself, which we approximate as being rigid. 6 Discrete to Continuum Modeling. 10 MIT, March, 2001 | Rosales. 2 1 d (cid:18) 1 n M L = M g sin (cid:18) n 2 N dt (cid:0) N (N + 1) (cid:20) (N + 1) (cid:20) + ((cid:18) (cid:18) ) ((cid:18) (cid:18) ) ; (3.6) +1 1
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) + ((cid:18) (cid:18) ) ((cid:18) (cid:18) ) ; (3.6) +1 1 n n n n(cid:0) ‘L ‘L (cid:0) (cid:0) (cid:0) for n = 2; : : : ; N 1; and (cid:0) 2 1 d (cid:18) 1 (N + 1) (cid:20) N M L = M g sin (cid:18) ((cid:18) (cid:18) ) : (3.7) 1 N N N (cid:0) 2 N dt N ‘L (cid:0) (cid:0) (cid:0) These are the equations for N torsion coupled equal pendulums. Remark 3.2 To check that the signs for the torsion forces selected in these equations are correct, take the di(cid:11)erence between the n and (n + 1) equation. Then you should see that the torsion th th force (due to the portion of the axle connecting the n and (n + 1) pendulums) is acting so as to th th make the angles equal. Remark 3.3 Note that the equations for the (cid:12)rst and last angle are di(cid:11)erent, because the (cid:12)rst and last pendulum experience a torsion force from only one side. How would you modify these equations to account for having one (or both) ends of the axle (cid:12)xed? Continuum Limit. Now we consider the continuum limit, in which we let N and assume that the n angle th can be written in the form: ! 1 where (cid:18) = (cid:18)(x; t) is a \nice" function (with derivatives) and x = ‘ is the position of the n N + 1 pendulum along the axle. In particular, note that: (cid:18) (t) = (cid:18)(x ; t) ; (3.8) n n n + 1 2 (cid:1)x = x x = : (3.9) +1 n n (cid:
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n + 1 2 (cid:1)x = x x = : (3.9) +1 n n (cid:0) N + 1 ‘ Take equation (3.6), and multiply it by N=‘. Then we obtain 2 d (cid:18) N (N + 1)(cid:20) n (cid:26) L = (cid:26) g sin (cid:18) + ((cid:18) 2(cid:18) + (cid:18) ) ; +1 1 n n n n(cid:0) 2 2 dt ‘ L (cid:0) (cid:0) where (cid:26) = M =‘ is the mass density per unit length in the N limit. Using equation (3.9), this can be written in the form: ! 1 2 d (cid:18) N (cid:20) (cid:18) 2(cid:18) + (cid:18) +1 1 n n n n(cid:0) (cid:26) L = (cid:26) g sin (cid:18) + : (3.10) n (cid:0) 2 2 dt (N + 1) L ((cid:1)x) (cid:0) From equation (3.8) we see that | in the limit N (where (cid:1) 0) | we have: ! 1 ! (cid:18) 2(cid:18) + (cid:18) @ (cid:18) +1 1 n n n(cid:0) 2 (cid:0) (x ; t) : n 2 2 ((cid:1)x) @x ! Discrete to Continuum Modeling. 11 MIT, March, 2001 | Rosales. Thus, (cid:12)nally, we obtain (for the continuum limit) the nonlinear wave equation (the \Sine{Gordon" equation): 2 2 (cid:18) c (cid:18) = ! sin (cid:18) ; (3.11) tt xx (cid:0) (cid:0) g (cid:20) where ! = is the pendulum angular frequency, and c = is a wave propagation speed 2 L
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(cid:0) g (cid:20) where ! = is the pendulum angular frequency, and c = is a wave propagation speed 2 L (cid:26)L s r (check that the dimensions are correct). Remark 3.4 Boundary Conditions. What happens with the (cid:12)rst (3.5) and last (3.7) equations in the limit N ? ! 1 As above, multiply (3.5) by 1=‘. Then the equation becomes: 2 (cid:26) L d (cid:18) (cid:26) g (N + 1)(cid:20) (cid:26) g (cid:20) (cid:18) (cid:18) 1 2 1 = sin (cid:18) + ((cid:18) (cid:18) ) = sin (cid:18) + : 1 2 1 1 (cid:0) 2 2 N dt N ‘ L N ‘L (cid:1)x (cid:0) (cid:0) (cid:0) Thus, as N one obtains ! 1 (cid:18) (0; t) = 0: x This is just the statement that there are no torsion forces at the x = 0 end (since the axle is free to rotate there). Similarly, one obtains: (cid:18) (‘; t) = 0; x at the other end of the axle. How would these boundary conditions be modi(cid:12)ed if the axle where (cid:12)xed at one (or both) ends? Kinks and Breathers for the Sine Gordon Equation. Equation (3.11), whose non-dimensional form is (cid:18) (cid:18) = sin (cid:18) ; (3.12) tt xx (cid:0) (cid:0) has a rather interesting history. Its (cid:12)rst appearance is not in the context of a physical context at all, but in the study of the geometry of surfaces with constant negative Gaussian curvature. Physical problems for which it has been used include: Josephson junction transmission lines, dislocation in crystals, propagation in ferromagnetic materials of waves carrying rotations in the magnetization direction, etc. Mathematically, it is a very interesting because it is one of the
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ferromagnetic materials of waves carrying rotations in the magnetization direction, etc. Mathematically, it is a very interesting because it is one of the few physically 10 important nonlinear partial di(cid:11)erential equations that can be solved explicitly (by a technique known as Inverse Scattering, which we will not describe here). An important consequence of equation (3.12) exact solvability, is that it possesses particle-like solutions, known as kinks, anti-kinks, and breathers. These are localized traveling distur- bances, which preserve their identity when they interact. In fact, the only e(cid:11)ect of an interaction is a phase shift in the particle positions after the interaction: e(cid:11)ectively, the "particles" approach each other, stay together brie(cid:13)y while they interact (this causes the "phase shift") and then depart, preserving their identities and original velocities. This can all be shown analytically, but here we will only illustrate the process, using some computational examples. 10 For reviews see: A. C. Scott, 1970, , Wiley Interscience, New York (page 250). Active and Nonlinear Wave Propagation in Electronics Barone, A. F. Esposito, C. J. Magee, and A. C. Scott, 1971, , Theory and Applications of the Sine Gordon Equation Rivista del Nuovo Cimento , pp. 227{267. vol. 1 Discrete to Continuum Modeling. 12 MIT, March, 2001 | Rosales. The (cid:12)rst step is to present analytical expressions for the various particle-like solutions of equation (3.12). These turn out to be relatively simple to write. Example 3.1 Kinks and Anti-Kinks. Equation (3.12) has some interesting solutions, that correspond to giving the pendulums a ful l 2(cid:25) twist (e.g.: take one end pendulum, and give it a ful l 2(cid:25) rotation). This generates a 2(cid:25) twist wave that propagates along the pendulum chain. These waves are known as kinks or anti-kinks (depending on the sign of the rotation), and can be written explicitly. In fact, they are steady wave solutions,
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as kinks or anti-kinks (depending on the sign of the rotation), and can be written explicitly. In fact, they are steady wave solutions, 11 for which the equation reduces to an O.D.E., which can be explicitly solved. Let 1 < c < 1 be a constant (kink, or anti-kink speed), and let z = (x c t x ) be a moving 0 (cid:0) (cid:0) (cid:0) coordinate, where the solution is steady | the "twist" wil l be centered at x = c t + x , where x is 0 0 the position at time t = 0. Then the kink solution is given by 2 z=(cid:12) e 1 z (cid:18) = 2 arccos = 4 arctan exp ; (3.13) 2 z=(cid:12) (cid:0) e + 1 (cid:0) (cid:12) ! !! where (cid:12) = 1 c is the kink width. This solution represents a propagating clock-wise 2(cid:25) rotation, p 2 (cid:0) from (cid:18) = 2 m (cid:25) as x (where m is an integer) to (cid:18) = 2 (m 1) (cid:25) as x , with most of the rotation concentrated in a region of width O((cid:12) ) near x = c t + x . The parameter c is determined 0 ! (cid:0)1 (cid:0) ! 1 (for example) by how fast the initial twist is introduced when the kink is generated. We note now that: 2c (cid:18) From (3.13) it fol lows that (cid:18) = c (cid:18) = sin . Using this, it is easy to show that (3.13) t x (cid:15) (cid:0) (cid:12) 2 ! is a solution of equation (3.12). The Sine-Gordon equation is the simplest of a "class" of models proposed for nuclear inter- (cid:15) actions. In this interpretation, the kinks are nuclear particles. Since (in the non-dimensional version (3.
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for nuclear inter- (cid:15) actions. In this interpretation, the kinks are nuclear particles. Since (in the non-dimensional version (3.12)) the speed of light is 1, the restriction 1 < c < 1 is the relativistic restriction, and the factor (cid:12) incorporates the usual relativistic contraction. (cid:0) The anti-kink solution fol lows by replacing x x and t t in (3.13). It corresponds to a propagating counter-clock-wise 2(cid:25) rotation, and it is given by ! (cid:0) ! (cid:0) 2 z=(cid:12) 1 e z (cid:18) = 2 arccos = 4 arctan exp : (3.14) (cid:0) 2 z=(cid:12) 1 + e (cid:12) ! !! The kinks and anti-kinks are very non-linear solutions. Thus, it is of some interest to study how they interact with each other. Because they are very localized solutions (non-trivial only in a smal l region), when their centers are far enough they can be added. Thus, numerical ly it is rather easy to study their interactions, by setting up initial conditions that correspond to kinks and anti-kinks far enough that they do not initial ly interact. Then they are fol lowed until they col lide. In the lectures the results of numerical experiments of this type wil l be shown (the numerical method used in the experiments is is a "pseudo-spectral" method). 11 (cid:0) Solutions of the form = ( ), where is a constant: the speed of propagation. (cid:18) (cid:18) x c t c Discrete to Continuum Modeling. 13 MIT, March, 2001 | Rosales. Example 3.2 Breathers. A di(cid:11)erent kind of interesting solution is provided by the "breathers" | which we hand le next. A breather is a wave-package kind of solution (an oscil latory wave, with an envelope that limits the wave to reside in a bounded region of space. These solutions vanish (exponential ly) as x . This last property al lows for easy numerical simulations of interactions of breathers (and kinks
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space. These solutions vanish (exponential ly) as x . This last property al lows for easy numerical simulations of interactions of breathers (and kinks). ! (cid:6)1 One can setup initial conditions corresponding to the interaction of as many kinks and/or breathers as one may wish (limited only be the numerical resolution of the computation), simply by separating them in space. A breather solution is characterized by two arbitrary constants 1 < d; V < 1. Then de(cid:12)ne A = d= 1 d ; p 2 B = 1= 1 V ; p (cid:0) 2 C = 1 d ; p (cid:0) 2 (cid:0) 9 > > > > (cid:0) > p = C B (V x t + t ) ; 0 q = d B (x V t x ) ; 0 (cid:0) (cid:0) (cid:0) Q = A sin(p)= cosh(q ) ; > > > > > > > > = > > > > > > > > > > > > > ; (3.15) where x and t are constants, centering the envelope and the phase, respectively. Notice that the 0 0 partial derivatives of Q (with respect to p and q) are given by Q = A cos(p)= cosh(q ) and Q = Q tanh(q ) : (3.16) p q (cid:0) The breather solution (and its time derivative) is then given by: (cid:18) = 4 arctan(Q) ; (cid:18) = 4 (1 + Q ) (C B Q + d B V Q ) : t p q 2 (cid:0) 9 > = > ; (3.17) The breather solution is a wave-package type of solution, with the phase control led by p, and the envelope (causing the exponential vanishing of the solution) by q). The wave-package details are given by: speed . . . . . . . . . . . . . . c = 1=V , p period . . . . . . . . . . . . . . T = 2 (cid:25)=(B C ) ; (3
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period . . . . . . . . . . . . . . T = 2 (cid:25)=(B C ) ; (3.18) p Phase. wave-length . . . . . . . . . (cid:21) = 2 (cid:25)=(B C V ) ; p 9 > = > ; speed . . . . . . . . . . . . . . c = V , e Envelope. (3.19) width . . . . . . . . . . . . . . e ) (cid:21) = 2 (cid:25)=(d B ) ; Notice that, while the phase moves faster than the speed of "light" (i.e.: 1), the envelope always moves with a speed 1 < V < 1, and has width proportional to 1 V . p 2 (cid:0) (cid:0) Final ly, in case you are familiar with the notion of group speed, notice that (for the linearized Sine- Gordon equation: (cid:18) (cid:18) + (cid:18) = 0) we have: (group speed) = 1/(phase speed) | which is exactly tt xx (cid:0) the relationship satis(cid:12)ed by c = V and c = 1=V for a breather. This is because, for x large, the e p breathers must satisfy the linearized equation. Thus the envelope must move at the group velocity j j corresponding to the oscil lations wave-length. Discrete to Continuum Modeling. 14 MIT, March, 2001 | Rosales. Remark 3.5 Pseudo-spectral Numerical Method for the Sine-Gordon Equation. Here we wil l give a rough idea of a numerical method that can be used to solve the Sine-Gordon equation. This remark wil l only make sense to you if you have some familiarity with Fourier Series for periodic functions. The basic idea in spectral methods is that the numerical di(cid:11)erentiation of a (smooth) periodic func- tions can be done much more eÆciently (and accurately) on the "Fourier Side" | since there it amounts to term by term multiplication of the n Fourier coe
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ly (and accurately) on the "Fourier Side" | since there it amounts to term by term multiplication of the n Fourier coeÆcient by in. On the other hand, th non-linear operations (such as calculating the square, point by point, of the solution) can be done eÆciently on the "Physical Side". Thus, in a numerical computation using a pseudo-spectral method, al l the operations involving taking derivatives are done using the Fourier Side, while al l the non-linear operations are done directly on the numerical solution. The back-and-forth calculation of Fourier Series and their inverses is carried by the FFT (Fast Fourier Transform) algorithm | which is a very eÆcient algorithm for doing Fourier calculations. Unfortunately, a naive implementation of a spectral scheme to solve the Sine-Gordon equation would require periodic in space, solutions. But we need to be able to solve for solutions that are mod-2(cid:25) periodic (such as the kinks and anti-kinks), since the solutions to the equation are angles. Thus, we need to get around this problem. In a naive implementation of a spectral method, we would write the equation as u = v ; t v = u sinu ; t xx ) (cid:0) (3.20) where u = (cid:18) and v = (cid:18) . Next we would discretize space using a periodic uniform mesh (with a large t enough period), and would evaluate the right hand side using FFT’s to calculate derivatives. This would reduce the P.D.E. to some large O.D.E., involving al l the values of the solution (and its time derivative) at the nodes in the space grid. This O.D.E. could then be solved using a standard O.D.E. solver | say, ode45 in MatLab. In order to use the idea above in a way that al lows us to solve the equation with mod-2(cid:25) periodicity in space, we need to be able to evaluate the derivative u in a way that ignores jumps by multiples xx of 2(cid:25) in u. The following trick works in doing this: Introduce U = e : Then iu gives a formula for u that ignores 2(cid:25) jumps in u. Warning: In the
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U = e : Then iu gives a formula for u that ignores 2(cid:25) jumps in u. Warning: In the actual implementation one xx 2 (U ) U U x xx u = i (3.21) xx (cid:0) 2 U must use 2 (U ) U U x xx u = imag xx (cid:0) 2 (cid:0) U ! to avoid smal l imaginary parts in the answer (caused by numerical errors). 4 Suggested problems. A list of suggested problems that go along with these notes follow: Discrete to Continuum Modeling. 15 MIT, March, 2001 | Rosales. 1. Check the derivation of the system of equations (2.20). 2. Derive the continuum equation in (2.21). 3. Look at the end of section 2, under the title "General Motion: String and Rods". Derive continuum equations describing the motion (in the plane) of a string without constraints. 4. Look at the end of section 2, under the title "General Motion: String and Rods". Add bending springs to the model, and derive continuum equations describing the motion (in the plane) of a rod without constraints. 5. Do the check stated in remark 3.2. 6. Answer the question in remark 3.3. 7. Do the dimensions check stated below equation (3.11). 8. Answer the question in remark 3.4. 9. Show that (3.13) is a solution (there is a hint about how to do this a few lines below the equation). 10. Use a computer to plot the solution in (3.13), as a function of z , for a few choices of c. 11. Show that (3.17) is a solution. 12. Use a computer to plot the solution in (3.17), as a function of x, for various times and choices of parameters. 13. Implement a numerical code to calculate interactions of kinks, breathers, etc., using the ideas sketched in remark 3.5.
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MIT OpenCourseWare http://ocw.mit.edu 8.512 Theory of Solids II Spring 2009 For information about citing these materials or our Terms of Use, visit: http://ocw.mit.edu/terms. Lecture 6: Scaling Theory of Localization 6.1 Notion of dimensionless conductance. g = G e2/π� G = σLd−2 where G is (resistance)−1 L = dimension of system Recall Einstein’s formula which relates conductivity to diffusion: where σ = 2e 2 N (0) Ω · D N (0) = # of states/energy Ω = unit volume D = diffusion coefficient Way to derive above equation is to calculate charge current density in two ways. j = −eD� · n j = −σ�V dV dn = −σ · �n · · · · · · current as diffusion of charge carriers V is local electrical potential Now, the local chemical potential is linearly related to V , i.e. µ = eV + constant ∴ = dn dv dn = e = e · dµ 2N (0) Ω ∴ j = − σ Ω e 2N (0) · �n Comparing this to Einstein’s formula, one gets · · · 2 is for spin. N (0) σ = 2e 2 Ω ⇒ G = e 2 2N (0) Ω · D · D · Ld−2 = 2e2 �D · � L2 · N (0) 1 (1)
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d−2 = 2e2 �D · � L2 · N (0) 1 (1) (2) 2 (3) 6.2 Thouless Energy The last equation suggests to define a quantity called Thouless Energy, ET ET = �D L2 = � τT N (0) ∼ 1 Δ · · · Δ is level spacing ⇒ G ∼ 2e2 ET � Δ ⇒ g ∼ ET Δ Now, Physical interpretation Assume the one box problem is solved and we want to study the behavior of the system as it gets bigger. Each box has its own distribution of energy levels. As the two boxes are in contact, the wavefunctions that were initially localized in separate boxes would mix. For weak coupling, one can do perturbation theory if mixing is smaller than typical level spacing Δ. For weak mixing, the wavefunctions resemble the initial states and so one gets localized distribution. For strong coupling, the final state would be a complicated mixture of all initial states and we get extended states. Physically one can think of diffusion as lifetime to leak particle out of box. And so it’s reasonable to think of Thouless energy as effective coupling. Thus the ratio ET /Δ, which also represents the dimensionless conductance ’g,’ characterizes the further scaling properties. 3 6.3 Sensitivity of energy eigenvalues to boundary conditions and its relation
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properties. 3 6.3 Sensitivity of energy eigenvalues to boundary conditions and its relation to ET . Here we derive another form of ET . Usually one uses periodic boundary conditions ψα(x + L) = ψα(x) · · · ψα is an eigenstate of the Hamiltonian. One can as well use the twisted boundary condition ψα(x + L) = e iφψα(x) · · · φ is same for all α. Such boundary conditions can arise if a magnetic field is present. Consider applying a uniform magnetic field (magnetic field is confined in flux tube) along the z direction. And we study the behavior of the eigenstates sitting on a cylindrical surface whose axis is parallel to the magnetic field. 1 � 2m i � + eA c 2 � ψα + V (r)ψα = Eαψα x ψ → ψ� exp −i A · d� � � o � ⇒ 1 2m −�2 ψ� α(r) + v(r)ψ� α(r) = Eαψ� α(r) Thus the gauge transformation has removed A from the Schrodinger equation, but now � � � ψα(r) has different boundary conditions. 4 ψ�(x + L) = ψ�(x) exp L i A · d� o � � = ψ�(x) exp [iφB] � where φB is magnetic �
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o � � = ψ�(x) exp [iφB] � where φB is magnetic flux. We will compute ΔEα using perturbation theory. φ L To first order in H �, ΔEα = 0 as < H � >∼< Vx >= 0. H = VxA = Vx e c � For perturbation theory second order in A, we have two terms: the diamagnetic term which gives a constant shift to all energy levels e2 A2 c2 2m The paramagnetic term gives, to second order in A (or φ), 1 1 2 m = φ2 | < β|VxA|α > |2 Eα − Eβ β±α � = φ2 L2 ∴ ∂2Eα ∂φ2 = 1 2 mL + 2 L2 | < β|Vx|α > |2 Eα − Eβ β�=α � | < β|Vx|α > |2 Eα − Eβ α�=β � (4) Thus one gets the same matrix element < β|Vx|α > that appeared in derivation of Kubo formula. Typical variation in energy levels can look like this. As first order shift is zero, all curves have zero slope at α = 0. All low lying states don’t undergo much variation as they occur in potential valleys and hence localized. One can do a sanity check on Eq.(4). On general grounds, there is equal probability for any level to go
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on Eq.(4). On general grounds, there is equal probability for any level to go up or down and so average variation in Eα’s should be zero. By using Thomas-Reich-Kunz F -sum rule, one can indeed show that ∂2Eα = 0 2∂φ α � 5 Now, | < β|Vx|α > |2 ∼ V 2 . Fluctuation of such terms is dominated by the term with maximum weight, i.e., by smallest Eα − Eβ variance � β�=α � | < β|Vx|α > |2 Eα − Eβ ∼ V 2 Δ � This is the average fluctuation in the curvature. Thus typically, ∼ Et ∂2Eα ∂φ2 � � � � Et ∼ 1 V 2 � L2 Δ � � � � 1 V 2 L2 Δ ∼ 1 2 V 2 · N(0) L Now, σ = 2e 2π N 2(0) Ω · V 2 = 2e 2 N(0) Ω D By comparing Boltzmann & Einstein formula for σ ⇒ D = πV 2N(0) ⇒ Et = D 2 L ∼ ET Thus for large ET /Δ, one gets extended states and hence large conductance, while for small E T /Δ, one gets localized states and hence small conductance. Thus we have established the relevance of Thouless energy for
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states and hence small conductance. Thus we have established the relevance of Thouless energy for conductance. 6.4 Application to 1 − Δ System 6 One expects Ohm’s low behavior for small length. g(L) ∼ 1 L as L increases, g(L) decreases as and so the states are localized. g(L) < 1 ⇒ ET Δ < 1 ⇒ For large L, g(L) ∼ e −L/ξ · · · ξ is the localization length Due to inelastic scattering the eigenvalues become ill-defined because of dephasing. Let us define a quantity τφ as the time over which a state looses the information about its phase. τ φ for various scattering mechanisms can be given by 1 τee ∼ T 2 �F 1 τe−phonon ∼ T 3 ω2 D Let us define Lφ as the length over which the dephasing occurs. So our argument works if Lφ = Doτφ � Lφ >> width of the sample. 7 Otherwise we are considering a system to which our derivation based on Quantum Mechanics ideas doesn’t apply and the system should rather be treated classically. And in such a situation, one used Ohm’s law. So to observe non-metallic behavior of copper wire, one needs to have a very small wire thickness or very low temperature as Lφ increases with decrease in temperature. Thus one expects the following behavior
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low temperature as Lφ increases with decrease in temperature. Thus one expects the following behavior of conductivity of a one-dimensional system. For Lφ > ξ, we can think of conductivity in terms of hopping of electrons over a distance scale ξ in time τφ. This gives rise to a diffusion constant of ξ2/τφ and resistivity which is proportional to τφ � T −p. [reference, book by Y. Imry]
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Lecture 3 Properties of MLE: consistency, asymptotic normality. Fisher information. In this section we will try to understand why MLEs are ’good’. Let us recall two facts from probability that we be used often throughout this course. • Law of Large Numbers (LLN): If the distribution of the i.i.d. sample X1, . . . , Xn is such that X1 has finite expectation, i.e. EX1 < , then the sample average | | → X1 + . . . + Xn n converges to its expectation in probability , which means that for any arbitrarily small α > 0, ¯ Xn = EX1 � ¯P( | Xn − EX1| > θ) 0 as n . � → � Note. Whenever we will use the LLN below we will simply say that the average converges to its expectation and will not mention in what sense. More mathematically inclined students are welcome to carry out these steps more rigorously, especially when we use LLN in combination with the Central Limit Theorem. Central Limit Theorem (CLT): If the distribution of the i.i.d. sample X1, . . . , Xn is such that X1 has finite expectation and variance, i.e. , then < EX1| | and π2 = Var(X) < → ≥n(X¯n − EX1) � → d N (0, π2) converges in distribution to normal distribution with zero mean and variance π2 , which means that for any interval [a, b], ≥n(X¯n − EX1) ∞ P � [a, b] � 16 � a � b 1 ≥2 ∂π e− 2 x 2�2 dx. • In other words, the random variable ≥n(X¯n − from normal distribution when n gets large. EX1) will behave like a random variable Exercise
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words, the random variable ≥n(X¯n − from normal distribution when n gets large. EX1) will behave like a random variable Exercise. Illustrate CLT by generating 100 Bernoulli random varibles B(p) (or one EX1). Repeat this many times Binomial r.v. B(100, p)) and then computing ≥n(X¯n − and use ’dfittool’ to see that this random quantity will be well approximated by normal distribution. We will prove that MLE satisfies (usually) the following two properties called consistency and asymptotic normality. 1. Consistency. We say that an estimate ϕˆ is consistent if ϕˆ ϕ0 in probability as , where ϕ0 is the ’true’ unknown parameter of the distribution of the sample. � n � → 2. Asymptotic Normality. We say that ϕˆ is asymptotically normal if ≥n(ϕˆ 2 d N(0, π� 0 ) 0 is called the asymptotic variance of the estimate ϕˆ. Asymptotic normality 2 where π� says that the estimator not only converges to the unknown parameter, but it converges fast enough, at a rate 1/≥n. ϕ0) � − Consistency of MLE. To make our discussion as simple as possible, let us assume that a likelihood function is smooth and behaves in a nice way like shown in figure 3.1, i.e. its maximum is achieved at a unique point ϕ. ˆ ) ϕ ( � 1 0.8 0.6 0.4 0.2 0 −0.2 −0.4 −0.6 −0.8 −1 �(ϕ) PSfrag replacements 0 0.5 1
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−0.6 −0.8 −1 �(ϕ) PSfrag replacements 0 0.5 1 1.5 2 ϕˆ 2.5 3 3.5 4 ϕ Figure 3.1: Maximum Likelihood Estimator (MLE) Suppose that the data X1, . . . , Xn is generated from a distribution with unknown pa­ rameter ϕ0 and ϕˆ is a MLE. Why ϕˆ converges to the unknown parameter ϕ0? This is not immediately obvious and in this section we will give a sketch of why this happens. 17 First of all, MLE ϕˆ is the maximizer of Ln(ϕ) = n1 n i=1 � log f (Xi| ϕ) 1 (of course, this does not affect maxi­ which is a log-likelihood function normalized by n mization). Notice that function Ln(ϕ) depends on data. Let us consider a function l(X ϕ) = log f (X ϕ) and define L(ϕ) = E�0 l(X ϕ), where E�0 denotes the expectation with respect to the true uknown parameter ϕ0 of the sample X1, . . . , Xn. If we deal with continuous distributions then | | | By law of large numbers, for any ϕ, L(ϕ) = � (log f (x ϕ))f (x ϕ0)dx. | | Ln(ϕ) � E�0 l(X ϕ) = L(ϕ). | Note that L(ϕ) does not depend on the sample, it only depends on ϕ. We will need the following Lemma. We have that for any ϕ, Moreover, the inequality is strict,
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need the following Lemma. We have that for any ϕ, Moreover, the inequality is strict, L(ϕ) < L(ϕ0), unless L(ϕ) ≡ L(ϕ0). P�0 (f (X ϕ) = f (X ϕ0)) = 1. | | which means that P� = P�0 . Proof. Let us consider the difference L(ϕ) − L(ϕ0) = E�0 (log f (X log f (X ϕ) | − ϕ0)) = E�0 log | Since log t t − ≡ 1, we can write X f ( X f ( ) . ) ϕ | ϕ | 0 E�0 log X f ( X f ( ϕ | ϕ | 0 ) ) ≡ ) ) − E�0 � f (x X f ( X f ( ϕ | ϕ | ϕ)dx | 0 = − � � 1 = x f ( x f ( ϕ | ϕ | 0 ) ) − 1 f (x ϕ0)dx | � f (x � � ϕ0)dx = 1 | − � 1 = 0. Both integrals are equal to 1 because we are integrating the probability density functions. This proves that L(ϕ) 0. The second statement of Lemma is also clear. L(ϕ0) − ≡ 18 We will use this Lemma to sketch the consistency of the MLE. Theorem: Under some regularity conditions on the family of distributions, MLE ϕˆ is consistent, i.e. ϕˆ ϕ0 as n . � → � The statement of this Theorem is
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consistent, i.e. ϕˆ ϕ0 as n . � → � The statement of this Theorem is not very precise but but rather than proving a rigorous mathematical statement our goal here is to illustrate the main idea. Mathematically inclined students are welcome to come up with some precise statement. Ln(ϕ) L(ϕ) PSfrag replacements ˆϕ ϕ0 ϕ Figure 3.2: Illustration to Theorem. Proof. We have the following facts: 1. ϕˆ is the maximizer of Ln(ϕ) (by definition). 2. ϕ0 is the maximizer of L(ϕ) (by Lemma). 3. This situation is illustrated in figure 3.2. Therefore, since two functions Ln and L are getting closer, the points of maximum should also get closer which exactly means that ϕˆ ϕ we have Ln(ϕ) L(ϕ) by LLN. � � ϕ0. � Asymptotic normality of MLE. Fisher information. We want to show the asymptotic normality of MLE, i.e. to show that ≥n(ϕˆ ϕ0) − � d N(0, π2 M LE ) for some π2 M LE and compute π2 First, we need to introduce the notion called Fisher Information. M LE . This asymptotic variance in some sense measures the quality of MLE. Let us recall that above we defined the function l(X ϕ) = log f (X ϕ). To simplify the notations we will denote by l�(X ϕ), l��(X ϕ), etc. the derivatives of l(X ϕ) with respect to ϕ. | | | | | Definition. (Fisher information.) Fisher information of a
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to ϕ. | | | | | Definition. (Fisher information.) Fisher information of a random variable X with distribution P�0 from the family I(ϕ0) = E�0 (l�(X P� : ϕ { ∞ ϕ0))2 | is defined by � } � E�0 �ϕ � � log f (X 19 2 . ϕ) | �=�0 � � � � Remark. Let us give a very informal interpretation of Fisher information. The derivative l�(X ϕ0) = (log f (X | | ϕ0))� = X f ( � X f ( ) 0 ) 0 ϕ | ϕ | can be interpreted as a measure of how quickly the distribution density or p.f. will change when we slightly change the parameter ϕ near ϕ0. When we square this and take expectation, i.e. average over X, we get an averaged version of this measure. So if Fisher information is large, this means that the distribution will change quickly when we move the parameter, so the distribution with parameter ϕ0 is ’quite different’ and ’can be well distinguished’ from the distributions with parameters not so close to ϕ0. This means that we should be able to estimate ϕ0 well based on the data. On the other hand, if Fisher information is small, this means that the distribution is ’very similar’ to distributions with parameter not so close to ϕ0 and, thus, more difficult to distinguish, so our estimation will be worse. We will see precisely
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��cult to distinguish, so our estimation will be worse. We will see precisely this behavior in Theorem below. Next lemma gives another often convenient way to compute Fisher information. Lemma. We have, E�0 l��(X ϕ0) | � �2 E�0 �ϕ2 log f (X ϕ0) = | I(ϕ0). − Proof. First of all, we have ϕ) = (log f (X l�(X | = ϕ))� | X f ( � X f ( ) ϕ | ϕ ) | and (log f (X = ϕ))�� | Also, since p.d.f. integrates to 1, X f ( �� X f ( ) ϕ | ϕ ) | − (f f 2 . ϕ X )) ( � | 2 ) ϕ X ( | if we take derivatives of this equation with respect to ϕ (and interchange derivative and integral, which can usually be done) we will get, f (x ϕ)dx = 1, | � � �ϕ f (x ϕ)dx = 0 and | �2 �ϕ2 f (x ϕ)dx = | � � � ϕ)dx = 0. f ��(x | To finish the proof we write the following computation E�0 l��(X ϕ0) = E�0 | �2 �ϕ2 x f ( �� x f ( � f ��(x | log f (X ϕ0) = | ) 0 ) ϕ | ϕ | ϕ0)dx 0 − � ) ϕ f (x � 0 | ϕ (x f ) | � E�0 (l�(X 0 − = = � � (log f
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(x f ) | � E�0 (l�(X 0 − = = � � (log f (x ϕ0))��f (x ϕ0)dx | | 2 f (x ϕ0)dx | � � ϕ0))2 = 0 | I(ϕ0 = I(ϕ0). − − 20 We are now ready to prove the main result of this section. Theorem. (Asymptotic normality of MLE.) We have, ≥n(ϕˆ ϕ0) − � N 0, � 1 I(ϕ0) � . As we can see, the asymptotic variance/dispersion of the estimate around true parameter will be smaller when Fisher information is larger. Proof. Since MLE ϕˆ is maximizer of Ln(ϕ) = n L� (ϕˆ) = 0. 1 n � Let us use the Mean Value Theorem n i =1 log f (Xi| ϕ), we have f (a) a f (b) b − − = f �(c) or f (a) = f (b) + f �(c)(a − b) for c [a, b] ∞ with f (ϕ) = L�n(ϕ), a = ϕˆ and b = ϕ0. Then we can write, 1)(ϕˆ 0 = L� (ϕˆ) = L�n(ϕ0) + L��(ϕˆ n n ϕ0) − for some ˆ ϕ1 ∞ [ˆ ϕ, ϕ0]. From here we get that ˆ ϕ − ϕ0 = 0) L� (ϕ n ˆ L��n ϕ1) ( − and ≥n(ϕˆ ϕ0) = − ≥nL (ϕ0) � n
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and ≥n(ϕˆ ϕ0) = − ≥nL (ϕ0) � n . ˆ Ln�� ϕ1) ( − (3.0.1) Since by Lemma in the previous section we know that ϕ0 is the maximizer of L(ϕ), we have L�(ϕ0) = E�0 l�(X ϕ0) = 0. (3.0.2) | Therefore, the numerator in (3.0.1) ≥nLn� (ϕ0) = ≥n � = ≥n 1 n 1 n n i=1 � n l�(Xi| ϕ0) − l�(Xi| ϕ0) − � E�0 l�(X1| � converges in distribution by Central Limit Theorem. i=1 � 0 (3.0.3) ϕ0) � � � N 0, Var�0 (l�(X1| ϕ0)) � Next, let us consider the denominator in (3.0.1). First of all, we have that for all ϕ, L��n(ϕ) = 1 n l��(Xi| ϕ) � E�0 l��(X1| ϕ) by LLN. (3.0.4) ϕ, ϕ0] and by consistency result of previous section, ˆ [ˆ ϕ ϕ0, we have ˆ ϕ1 � � ϕ0. � 1 ∞ Also, since ϕˆ Using this together with (10.0.3) we get E�0 l��(X1| L��n(ϕˆ 1) � ϕ0) = I(ϕ0) by Lemma above. − 21 Combining this with (3.0.3) we get ≥nL�n(ϕ0) − L��n(ϕˆ1
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3.0.3) we get ≥nL�n(ϕ0) − L��n(ϕˆ1) � d N 0, � Var�0 (l�(X1| (I(ϕ0))2 ϕ0)) . � Finally, the variance, Var�0 (l�(X1| ϕ0)) = E�0 (l�(X ϕ0))2 | (E�0 l�(x ϕ0))2 = I(ϕ0) | − − 0 where in the last equality we used the definition of Fisher information and (3.0.2). Let us compute Fisher information for some particular distributions. Example 1. The family of Bernoulli distributions B(p) has p.f. f (x p) = p x(1 | − p)1 x − and taking the logarithm The second derivative with respect to parameter p is log f (x p) = x log p + (1 | − x) log(1 − p). � �p log f (x p) = | x p − 1 1 �2 x p �p2 , log f (x p) = | − x 2 p − 1 (1 x . p)2 − − Then the Fisher information can be computed as I(p) = E − �2 �p2 ¯ log f (X p) = | EX p2 + EX p)2 1 − (1 − = p p2 + 1 (1 p p)2 − − 1 p(1 − . p) The MLE of p is ˆp = X and the asymptotic normality result states that − − = which, of course, also follows directly from the CLT. Example. The family of exponential distributions E(�) has p.d.f. ≥n(ˆp − p0) � N(0, p0(1
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) has p.d.f. ≥n(ˆp − p0) � N(0, p0(1 − p0)) f (x �) = | � �e− 0, �x , x 0 ∀ x < 0 and, therefore, − This does not depend on X and we get log f (x �) = log � | �x ≤ �2 ��2 log f (x | �) = 1 − �2 . 1 . �2 Therefore, the MLE ˆ� = 1/X¯ is asymptotically normal and �) = | log f (X I(�) = �2 ��2 E − ≥n(ˆ� − �0) � N(0, �2).0 22
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8.701 0. Introduction 0.6 Particles Introduction to Nuclear and Particle Physics Markus Klute - MIT 1 Introduction 10-10m a few 10-15m 10-15m 1fm 2 Force Particles © Mattson Rosenbaum (on PBworks). All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. 3 Matter Particles © Mattson Rosenbaum (on PBworks). All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. 4 The Higgs Boson © CERN. All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. © Nature Publishing Group. All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. 5 Elementary Particle © CERN. All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. 6 Timeline of Discoveries © The Economist Newspaper Limited. All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. 7 Composite Particles and Hadrons Mesons: quark-antiquark states; bosons Pion Courtesy of Arpad Horvath on Wikimedia. License: CC BY- SA. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. © CERN. All rights reserved. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. Baryons: three-quark states; fermions These images are in the public domain. 8 Proton Courtesy of Arpad Horvath on Wikimedia. License: CC BY- SA. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. Nuclei Bound state of protons
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from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. Nuclei Bound state of protons and neutrons through the strong force. Can be described by number of protons, Z, (atomic number) and number of neutrons, N. The sum Z+N is denoted atomic mass A Courtesy of Sjlegg on Wikimedia. License: CC BY-SA. This content is excluded from our Creative Commons license. For more information, see https://ocw.mit.edu/fairuse. 9 MIT OpenCourseWare https://ocw.mit.edu 8.701 Introduction to Nuclear and Particle Physics Fall 2020 For information about citing these materials or our Terms of Use, visit: https://ocw.mit.edu/terms.
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Turbulent Flow and Transport 1 Review of Fundamental Laws and Constitutive Equations 1.1 Fundamental laws governing continuum flow, expressed in terms of (i) material volumes (closed systems) and (ii) control volumes. 1.2 Mass conservation equation; integral and differential forms. The equation of motion in terms of the stress tensor. Integral and differential 1.3 forms. Stress tensor for a Newtonian fluid. The Navier−Stokes equation. 1.4 The energy equation (First Law) in integral and differential forms. The equations for the kinetic, potential, and internal energies. Physical significance of the various terms. The viscous dissipation function. 1.5 The Second Law of thermodynamics. 1.6 1.7 1.8 1.9 The equation for entropy production (Gibbs’ equation). The Second Law as a statement that the viscous dissipation function is greater than or equal to zero. The thermodynamic equations of state: differential expressions for all the thermodynamic properties of a fluid in terms of three measurable properties which are functions of pressure and temperature−the coefficient of thermal expansion, the isothermal compressibility, and the specific heat at constant pressure. The differential equation for temperature and heat flux. Example: thermal effects due to viscous heating in Couette flow. Some low−speed approximations: (i) "incompressible flow," (ii) the neglect of the isentropic compression term in the temperature equation. Criteria for validity. 1.10 The conservation equation for a molecular species. Mass transfer. 1.11 Introduction to the molecular basis of viscosity, heat conduction and diffusion in terms of a simplified kinetic theory of gases. (This serves as an important but relatively simple analogy for the random−walk transport that also occurs in turbulence.) References: Sec. 1.1−1.4: Sonin. Fundamental Laws of Motion: Particles, Material Volumes, And Control Volumes . Sonin. The Equation of Motion for Viscous Fluids. available at http://web.mit.edu/2.25/www/ ; White.Viscous Fluid Flow , 2nd ed. (1991): 59−89, 96−100; Pope. Ch. 2; or other books. Secs 1
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nd ed. (1991): 59−89, 96−100; Pope. Ch. 2; or other books. Secs 1.4 − 1.6: Class notes plus summaries handed out by Sonin. Sec. 1.10: Sec. 1.11: For Sec. 1.7−1.8: Handout: Sonin."The Thermodynamic Constitutive Equations and the Equation for Temperature." Class notes. See for example Bird, Curtis, and Hirschfelder. Molecular Theory of Gases and Liquid, (1954): 8−16,orVincenti&Kruger. Introduction to Physical Gas Dynamics. 1965:15−20.
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6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. 6.001 Notes: Section 7.1 Slide 7.1.1 In the past few lectures, we have seen a series of tools for helping us create procedures to compute a variety of computational processes. Before we move on to more complex issues in computation, it is useful to step back and look at more general issues in the process of creating procedures. In particular, we want to spend a little bit of time talking about good programming practices. This sounds a little bit like lecturing about "motherhood and apple pie", that is, a bit like talking about things that seem obvious, apparent, and boring in that everyone understands and accepts them. However, it is surprising how many "experienced" programmers don't execute good programming practices, and we want to get you started on the right track. Slide 7.1.2 Thus, in this lecture we are going to look briefly at several methodological aspects of creating procedures: designing the components of our code, debugging our code when it doesn't run correctly, writing documentation for our code, and testing our code. We will highlight some standard practices for each stage, and indicate why these practices lead to efficient and effective generation of code. Slide 7.1.3 Let’s start with the issue of how to design code, given a problem statement. There are many ways to do this, but most of them involve some combination of the following steps: • • • Design of data structures Design of computational modules Design of interfaces between modules Once we have laid out the general design of these stages, we follow by creating specific instantiations of the actual components. We have not yet talked about data structures in Scheme, and will return to this issue in a few lectures. For our purposes here, the key thing to note is that when designing a computational system, it is extremely valuable to decide what kinds of information naturally should be grouped together, and to then create structures that perform that grouping, while maintaining interfaces to the structures that hide the details. For example, one thinks naturally of a vector as a pairing of an x and y coordinate. One wants to 6.001 Structure and Interpretation of Computer Programs.
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thinks naturally of a vector as a pairing of an x and y coordinate. One wants to 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. be able to get out the coordinates when needed, but in many cases, one thinks naturally of manipulating a vector as a unit. Similarly, one can imagine aggregating together a set of vectors, to form a polygon, and again one can think of manipulating the polygon as a unit. Thus, a key stage in designing a computational system is determining the natural data structures of the system. Slide 7.1.4 A second stage in designing a computational system is deciding how best to break the computation into modules or pieces. This is often as much art as science, but there are some general guidelines that help us separate out modules in our design. For example, is there part of the problem that defines a computation that is likely to be used many times? Are there parts of the problem that can be conceptualized in terms of their behavior, e.g. how they convert certain inputs into certain types of outputs, without worrying about the details of how that is done. Does this help us focus on other parts of the computation? Or said a bit differently, can one identify parts of the computation in terms of their role, and think about that role in the overall computation, without having to know details of the computation? If one can, these parts of the computation are good candidates for separate modules, since we can focus on their use while ignoring the details of how they achieve that computation. Slide 7.1.5 Finally, given that one can identify data structures, whose information is to be manipulated; and stages of computation, in which that information is transformed; one wants to decide the overall flow of information between the modules. What types of inputs does each module need? What types of data does each module return? How does one ensure that the correct types are provided, in the correct order? These kinds of questions need to be addressed in designing the overall flow between the computational modules. Slide 7.1.6 This is perhaps more easily seen by thinking about an example – and in fact you have already seen one such example, our implementation ofsqrt. When we implemented our method for square roots, we actually engaged
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– and in fact you have already seen one such example, our implementation ofsqrt. When we implemented our method for square roots, we actually engaged in many of these stages. We didn’t worry about data structures, since we were simply interested in numbers. We did, however, spend some effort in separating out modules. Remember our basic computation: we start with a guess; if it is good enough, we stop; otherwise we make a new guess by averaging the current guess, and the ratio of the target number and the guess, and continue. To design this system, we separated out several modules: the notion of averaging, the notion of measuring “good enough”. We saw that some of these modules might themselves rely on other procedural abstractions; for example, 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. our particular version of “good enough” needed to use the absolute value procedure, though other versions might not. Slide 7.1.7 Once we had separated out these notions of different computations: average and good-enough, we considered the overall flow of information through the modules. Note by the way that we can consider each of theses processes as a black box abstraction, meaning that we can focus on using these procedures without having to have already designed the specific implementation of each. Now what about the flow between these modules? In our case, we began with a guess, and tested to see if it was good enough. If it was, we could then stop, and just return the value of the guess. Slide 7.1.8 If it was not, then we needed to average the current guess and the ratio of our target number to the guess. Slide 7.1.9 And then we need to repeat the entire process, with this new value as our new guess. The point of laying out these modules, or black boxes, is that we can use them to decide how to divide up the code, and how to isolate details of a procedure from its use. As we saw when we implemented our sqrt procedure, we can change details of a procedure, such as average, without having to change any of the procedures that use that particular component. As well, the flow of
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of a procedure, such as average, without having to change any of the procedures that use that particular component. As well, the flow of information between the modules helps guide us in the creation of the overall set of procedures. Thus, when faced with any new computational problem, we want to try to engage in the same exercise: block out chunks of the computation that can be easily isolated; identify the inputs and outputs from each chunk; and lay out the overall flow of information through the system. Then we can turn to implementing each of the units separately, and testing the entire system while isolating the effects of each unit. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.1.10 A second key element to good programming practice is code documentation. Unfortunately, this is one of the least well- practiced elements – far too often programmers are in such a hurry to get things written that they skip by the documentation stage. While this may seem reasonable at the time of code creation, when the design choices are fresh in the program creator’s mind, six months later when one is trying to read the code (even one’s own), it may be very difficult to reconstruct why certain choices were made. Indeed, in many commercial programming settings, more time is spent on code maintenance and modification than on code generation, yet without good documentation it can be very difficult or inefficient to understand existing code and change it. As well, good documentation can serve as a valuable source of information about the behavior of each module, enabling a programmer to maintain the isolation of the details of the procedural abstraction from the use of that abstraction. This information can be of help when debugging procedures. Slide 7.1.11 As with designing procedural modules, the creation of good documentation is as much art as science. Nonetheless, here are some standard elements of well-documented code. We are going to illustrate each of these with an example. Slide 7.1.12 First, describe the goal of the procedure. Is it intended to part of some other computation (as this helper function is)? If so, what is the rough description of the process? Note that here we have been a bit cryptic (in order to fit things on the slide
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, what is the rough description of the process? Note that here we have been a bit cryptic (in order to fit things on the slide) and we might well want to say more about “successive refinement” (though we could defer that to the documentation under the improve procedure). We also identify the role of each argument to the procedure. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.1.13 Second, describe the types of values used in the computation. In this case, the inputs or parameters are both numbers, and the returned value is also a number. Actually, if we were more careful here, we would require that X be a positive number, and we would place a check somewhere to ensure that this is true. Slide 7.1.14 Third, describe constraints, either desired or required, on the computation. Here, we know that squaring the guess should get us something close to the target value, although we really don’t guarantee this until we reach the termination stage. Slide 7.1.15 And fourth, describe the expected state of the computation and the goal at each stage in the process. For example, here we indicate what good-enuf? should do, namely test if our approximation is sufficiently accurate. Then we indicate that if this is the case, we can stop and what value to return to satisfy the contract of the entire procedure. And we indicate how to continue the process, though we could probably say a bit more about what improve should do. Notice how we can use the documentation to check some aspects of our procedure’s “contract”. Here, we have indicated that the procedure should return a number. By examining the if expression, we can see that in the consequent clause, if the input parameter guess is a number, then we are guaranteed to return a number. For the alternative clause, we can use induction to reason that given numbers as input, we also return a number, and hence the entire procedure returns a value of the correct type. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.1.16 In general, taking care
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and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.1.16 In general, taking care to meet each of the stages when you create code will often ensure an easier time when you have to refine or replace code. Getting into the habit of doing this every time you write something, even if you are only minutes away from some problem set deadline, will greatly improve your productivity! 6.001 Notes: Section 7.2 Slide 7.2.1 While we would like to believe that the code we write will always run correctly, the first time we try it, experience shows that this is a fortunate happenstance. Typically, especially with complex code, things will not work right, and we need to debug our code. Debugging is in part an acquired skill – with lots of practice you will develop your own preferred approach. Here, we are going to describe some of the common sources of errors in code, and standard tools for finding the causes of the errors and fixing them. Slide 7.2.2 A common and simple bug in code arises when we use an unbound variable. From the perspective of Scheme, this means that somewhere in our code we try to reference (or look up the value of) a variable that does not have one. This can occur for several reasons. The simplest is that we mistyped – a spelling error. The solution in this case is pretty straightforward – simply search through the code file using editor tools to find the offending instance and correct it. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.2.3 Sometimes, however, we are using a legal variable (that is, one that we intended to hold some value) but the evaluator still complains that this variable is unbound. How can that be? Remember that in Scheme a variable gets bound to a value in one of several ways. We may define it at “top level”, that is, we may directly tell the interpreter to give a variable some value. We may define it internally within some procedure. Or, we may use it as a formal parameter of a procedure, in which case it gets locally bound to a value when the
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some procedure. Or, we may use it as a formal parameter of a procedure, in which case it gets locally bound to a value when the procedure is applied. In the last two cases, if we attempt to reference the variable outside the scope of the binding, that is, somewhere outside the bounds of the lambda expression in which the variable is being used, we will get an unbound variable error. This means that we have tried to use a variable outside its legal domain, and we need to correct this. This probably means we have a coding error, but we can isolate the problem either by searching for instances of the variable in the code file, or by using the debugger. Slide 7.2.4 So what does a debugger do to help us find errors? Each programming language will have its own flavor of debugger; for an interpreted language like Scheme, the debugger actually places us inside the state of the computation. That is, when an error occurs, the debugger provides us access to the state of the computation at the time of the error, including access to the values of the variables within the computation. Moreover, we can step around inside the environment of the computation: we can work back up the chain of computational steps, examining what values were produced during reductions (where computation is reduced to a simpler expression), and examining what values were produced during substitutions (where the computation was converted to a simpler version of itself). Slide 7.2.5 For example, here is a simple procedure, which we have called with argument 2. Notice what happens when we hit the unbound variable error and enter the debugger. We are placed at the spot in the computation at which the error occurred. If we choose to step back through the chain of evaluations, we can see what expressions were reduced to get to this point, and what recursive versions of the same problem were invoked in reaching this stage. In this case, we note that foo was initially called with argument 2, and after a reduction through an if expression, we arrived at an expression that contained within it a simpler version of the same problem. This reduction stage repeated again, until we apparently reached the base case of the if expression, where we hit the unbound variable. We can see in this simple case that our unbound
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we apparently reached the base case of the if expression, where we hit the unbound variable. We can see in this simple case that our unbound error is coming from within the body of foo and is in the base case of the decision process. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.2.6 A second class of errors deals with mistakes in syntax – creating expressions that do not satisfy the programming language’s rules for creating legal expressions. A simple one of these is an expression in which the wrong number of arguments is provided to the procedure. If this occurs while attempting to evaluate the offending expression, we will usually be thrown into the debugger – a system intended to help us determine the source of the error. In Scheme, the debugger provides us with information about the environment in which the offending expression occurred. It supplies tools for examining the values associated with variable names, and for examining the sequence of expressions that have been evaluated leading up to this error. By stepping through the frames of the debugger, we can often isolate where in our code the incorrect expression resides. Slide 7.2.7 A more insidious syntax error occurs when we use an expression of the wrong type somewhere in our code. If we use an expression whose value is not a procedure as the first subexpression of a combination, we will get an error that indicates we have tried to apply a non-procedure object. As before, the debugger can often help us isolate the location of this error, though it may not provide much insight into why an incorrect object was used as a procedure. For that, we may have to trace back through our code, to determine how this value was supplied to the offending expression. The harder error to isolate is one in which one of the argument expressions to a combination is of the wrong type. The reason this is harder to track down is that the cause of the creation of an incorrect object type may have occurred far upstream, that is, some other part of our code may have created an incorrect object, which has been passed through several levels of procedure calls before causing an error. Tracking down the original source of this error can be difficult, as we need to chase our way back through
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of procedure calls before causing an error. Tracking down the original source of this error can be difficult, as we need to chase our way back through the sequence of expression evaluations to find where we accidentally created the wrong type of argument. Slide 7.2.8 The most common sorts of errors, though, are structural ones. This means that our code is syntactically valid – composed of correctly phrased expressions, but the code does not compute what we intended, because we have made an error somewhere in the code design. This could be for a variety of reasons: we started a recursive process with the wrong initial values, or we are ending at the wrong place, or we are updating parameters incorrectly, or we are using the wrong procedure somewhere, and so on. Finding these errors is tougher, since the code may run without causing a language error, but the results we get are erroneous. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.2.9 This is where having good test cases is important. For example, when testing a recursive procedure, it is valuable to try it using the base case values of the parameters, to ensure that the procedure is terminating at the right place, and returning the right value. It is also valuable to select input parameter values that sample or span the range of legal values – does it work with small values, with large values; does changing the input value by a small increment cause the expected change in output value? Slide 7.2.10 And what do we do if we find we have one of these structure errors? Well, our goal is to isolate the location of our misconception within the code, and to do this, there are two standard tools. The most common one is to use a print or display expression – that is, to insert into our code, expressions that will print out for us useful information at different stages of the computation. For example, we might insert a display expression within the recursive loop of a procedure, which will print out information about the values of parameters. This will allow us to check that parameters are being updated correctly, and that end cases are correctly seeking the right termination point. We might similarly print
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us to check that parameters are being updated correctly, and that end cases are correctly seeking the right termination point. We might similarly print out the values of intermediate computations within recursive loops, again to ascertain that the computation is operating with the values we expect, and is computing the values we expect. A related tool, supplied for example with Scheme, is a tracer. This allows us to ask the evaluator to inform us about the calling conventions of procedures – that is, to print out the values of the parameters supplied before each application of the procedure we designate, and the value returned by each such procedure call. This is similar to our use of display expressions, but is handled automatically for us. It applies only to parameters of procedure calls, however, so that if we want to examine for detailed states of the computation, we need to fall back on the display tactic. In some cases, it may help to actually walk through the substitution model, that is, to see each step of the evaluation. Many languages, including Scheme, provide a means for doing this – in our case called the stepper. This is a mechanism that lets us control each step of the substitution model in the evaluation of the expression. It is obviously tedious, but works best when we need to isolate a very specific spot at which an error is occurring, and we don’t want to insert a ton of display expressions. Perhaps the best way to see the role of these tools is to look at an example, which we do next. 6.001 Notes: Section 7.3 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.3.1 Let’s use an example of a debugging session to highlight these ideas. This will be primarily to fix a structural error, but we will see how the other tools come into play as we do this. Suppose we want to compute an approximation to the sine function. Here is a mathematical approximation that will give us a pretty good solution. So let’s try coding this up. Slide 7.3.2 So here is a first attempt at some code to do this. We will assume that fact and small-enuf? already exist. The basic idea behind this procedure is quite similar to what we did for square roots.
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that fact and small-enuf? already exist. The basic idea behind this procedure is quite similar to what we did for square roots. We start with a guess. We then see how to improve the guess, in this case by computing the next term in the approximation, which we would like to add in. If this improvement is small enough, we are done and can return the desired value. If not, we repeat the process with a better guess, by adding in the improvement to the current guess. Slide 7.3.3 Now, let’s try it out on some test cases. One nice test case is the base case, of x equal to 0. That clearly works. Another nice test case is when x is equal to pi, where we know the result should also be close to 0. Oops! That didn’t work. Nor does the code work for x equal to pi half. Both of these latter cases give results that are much too large. Slide 7.3.4 Okay, we need to figure out where our conceptual error lies. Let’s try to isolate this by tracing through the computation. In particular, we will add some display expressions that will show us the state of the computation each time through the recursion. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.3.5 And let’s try this again. Here we have used the test case of x equal to pi. And we can see the trace of the computation. If we compare this to the mathematical equation we can see one problem. We really only want terms where n is odd, but clearly we are getting all terms for n. So we need to fix this. Most likely this is because we are not changing our parameters properly. Slide 7.3.6 So here is the correction. We will need to increment our parameter by 2 each time, not by 1 – an easy mistake to make, and to miss! Slide 7.3.7 So let’s try this again. Hmm. We have gotten better as we are only computing the odd terms for n, but we are still not right. If we look again at
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. We have gotten better as we are only computing the odd terms for n, but we are still not right. If we look again at the mathematical equation, we can see that we should be alternating signs on each term. Or said another way, the successive approximations should go up, then down, then up, then down, and so on. Note that we could have also spotted this if we had chosen to display the value of next at each step. So we need to keep track of some additional information, in this case whether the term should be added or subtracted from the current guess. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.3.8 Well, we can handle that. We add another parameter to our helper procedure, which keeps track of whether to add the term (if the value is 1) or whether to subtract the term (if the value is ­ 1). And of course we will need to change how we update the guess, and how we update the value of this parameter. Slide 7.3.9 Oops! We blew it somewhere! We could enter the debugger to locate the problem, but we can already guess that since we changed the aux procedure, that must be the cause. Slide 7.3.10 And clearly the solution is to make sure we call this procedure with the right number of arguments. Notice that in this case it is easy to spot this error, but in general, we should get into the habit of checking all calls to a procedure when we alter its set of parameters. Slide 7.3.11 Now, if we try this on the test case of x equal pi, this works! But if we try it on the test case of pi half, it doesn’t! The answer should be close to 1, but we are getting something close to -1. Note that this reinforces why we want to try a range of test cases – if we had stopped with x equal pi, we would not have spotted this problem. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.3.12 Here is the bug. We started with the wrong initial value –
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2004 by Massachusetts Institute of Technology. Slide 7.3.12 Here is the bug. We started with the wrong initial value – a common error. By fixing this, we can try again and … Slide 7.3.13 … finally we get correct performance. Note how we have used printing of values to isolate changes, as well as using the debugger to find syntax errors. Slide 7.3.14 In general, we want you to get into the habit of doing the same things. Developing good programming methodology habits now will greatly help you when you have to deal with large, complex, bodies of code. Good programming discipline means being careful and thorough in the creation and refinement of code of all sizes and forms, so start exercising your “programming muscles” now! 6.001 Notes: Section 7.4 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.4.1 One other tool that we have in our armamentarium of debugging is the use of types. In particular, type specifications, that is, constraints on what types of objects are passed as arguments to procedures and what types of objects are returned as values by procedures, can help us both in planning and designing code, and in debugging existing code. Here, we are going to briefly explore both of these ideas, both to demonstrate why careful program practice can lead to efficient generation of robust code; and to illustrate why thinking about types of procedures and objects is a valuable practice. Slide 7.4.2 To motivate the idea of types as a tool in designing code, let's consider an example. Suppose we want to create a procedure, let's call it repeated, that will apply any other procedure some specified number of times. Since this is a vague description, let's look at a specific motivating example. We saw earlier the idea that we could implement multiplication as a successive set of additions, and that we could implement exponentiation as a successive set of multiplications. If we look at these two procedures, we can see that there is a general pattern here. There is a base case value to return (0 in one case, 1 in the other). And there is the idea of applying an operation to an input value
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a base case value to return (0 in one case, 1 in the other). And there is the idea of applying an operation to an input value and the result of repeating that process one fewer times. Repeated is intended to capture that common pattern of operation. Slide 7.4.3 So here is what we envision: we want our repeated procedure to take a procedure to repeat, and the number of times to repeat it. It should return a procedure that will actually do that, when applied to some value. Here we can see that the procedure being applied would change in each case, and the initial value to which to apply it would change, but otherwise the overall operation is the same. The question is: how to we create repeated, and why does the call to repeated have that funny structure, with two open parens? 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.4.4 First, what is the type of repeated? Well, from the previous slide, we know that as given, it should take two arguments. The first should be a procedure of one argument. We don't necessarily know what type of argument this procedure should take (in the two examples shown, the input was a number, but we might want to be more general than this). What we do know, is that whatever type of argument this procedure takes, it needs to return a value of the same type, since it is going to apply that procedure again to that value. Hence the first argument to repeated must be a procedure of type A to A. The second argument to repeated must be an integer, since we can only apply an operation an integer number of times. Actually, it should probably be a non-negative integer. And as we argued, the returned object needs to be a procedure of the same type: A to A because the idea is to use repeated recursively on itself. Slide 7.4.5 Okay, now how does this help us in designing the actual procedure? We know the rough form that repeated should take. It should have a test for the base case, which is when there are no more repetitions to make. In the base case, it needs to do something, which we have to figure out. And in the recursive case, we
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more repetitions to make. In the base case, it needs to do something, which we have to figure out. And in the recursive case, we expect to use repeated to solve the smaller problem of repetition, plus some additional operations, which we need to figure out. Slide 7.4.6 For the base case, what do we know? We know that by the type information, this must return a procedure of a single argument that returns a value of the same type. We also know that if we are in the base case, there is really nothing to do. We don't want to apply our procedure any more times. Hence, we can deduce that we need to return a procedure that serves as the identity: it simply returns whatever value was passed in. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.4.7 Now, what about the recursive case? Well, the idea is to apply the input procedure to the result of repeating the operation n-1 times. How do we use this idea to figure out the correct code? First, we know that whatever we write must have type A to A by the specification of repeated. Slide 7.4.8 Next we know that we want to apply the input procedure proc to the result of solving the same problem, n-1 times. So we ought to have something that has these pieces in it. Slide 7.4.9 But let's check the types. We know that repeated has type A to A, and the proc expects only an argument of type A. So clearly we need to apply repeated to an argument before passing the result on to proc. Hence we have the form shown. Note how this fairly complex piece of code can be easily deduced by using types of procedures to determine interactions. Of course, to be sure we did it right, we should now test this on some test cases, for example, by running mul or exp on known cases. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.4.10 A second way that types can help us is in debugging code. In particular, we can use the information about types of arguments and types of
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A second way that types can help us is in debugging code. In particular, we can use the information about types of arguments and types of return values explicitly to check that procedures are interacting correctly. And in some cases, where there are constraints on the actual values being returned, we can also enforce a check. Slide 7.4.11 As an example, here is our sqrt code from before. One of the conditions we have is that the input arguments need to be numbers. And we could check that numbers are being correctly passed in by inserting an explicit check. In this case, it is probably redundant since the only code that calls sqrt­ helper is itself, but in general, when multiple procedures might be involved, you can see how this check is valuable. Note that one can insert this check only when debugging, as a tool for deducing what procedure is incorrectly supplying arguments. But one can also use it regularly, if you want to ensure robust operation of the code. Clearly one could add a check on the return value in a similar fashion. Slide 7.4.12 But there are other things one could use to ensure correct operation. For example, the number whose square root we are seeking should be a positive number, and we could check that as shown. Thus we see that types also serve as a useful tool on good programming methodologies. 6.001 Structure and Interpretation of Computer Programs. Copyright © 2004 by Massachusetts Institute of Technology. Slide 7.4.13 To summarize, we have seen a set of tools for good programming practices: ways of designing code, debugging code, evaluating code, and using knowledge of code structure to guide the design.
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MIT OpenCourseWare http://ocw.mit.edu 6.641 Electromagnetic Fields, Forces, and Motion, Spring 2005 Please use the following citation format: Markus Zahn, 6.641 Electromagnetic Fields, Forces, and Motion, Spring 2005. (Massachusetts Institute of Technology: MIT OpenCourseWare). http://ocw.mit.edu (accessed MM DD, YYYY). License: Creative Commons Attribution-Noncommercial-Share Alike. Note: Please use the actual date you accessed this material in your citation. For more information about citing these materials or our Terms of Use, visit: http://ocw.mit.edu/terms 6.641, Electromagnetic Fields, Forces, and Motion Prof. Markus Zahn Lecture 1: Integral Form of Maxwell’s Equations I. Maxwell’s Equations in Integral Form in Free Space 1. Faraday’s Law d ∫(cid:118) E ds = - dt ∫ µ0 H i da i C S Circulation of E Magnetic Flux 0 µ = 4π ×10-7 henries/meter [magnetic permeability of free space] (Kirchoff’s Voltage Law, conservative electric EQS form: E ds = 0 i(cid:118) ∫ C field) MQS circuit form: v = L (Inductor) di dt 2. Ampère’s Law (with displacement current) (cid:118) ∫ H ds i = C i ∫ J da + S d dt ε 0 ∫ S i E da Circulation Conduction Displacement of H Current Current MQS form: H ds = ∫ J da i (cid:118)∫ i C EQS circuit form: i = C (capacitor) S dv dt 6.641, Electromagnetic Fields, Forces, and Motion Prof. Markus Zahn Lecture 1 Page 1 of 6 3. Gauss’ Law for Electric Field (cid:118)∫ ε
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Lecture 1 Page 1 of 6 3. Gauss’ Law for Electric Field (cid:118)∫ ε 0E da = ρ dV ∫ i S V ≈ 8.854 ×10-12 farads/meter ε 0 ≈ 10-9 36π 1 ε µ 0 0 free space) c = 3≈ × 108 meters/second (Speed of electromagnetic waves in 4. Gauss’ Law for Magnetic Field (cid:118)∫ µ 0H da = 0 i S In free space: B = µ H0 magnetic flux density (Teslas) magnetic field intensity (amperes/meter) 5. Conservation of Charge Take Ampère’s Law with displacement current and let contour C → 0 lim (cid:118)∫ H i ds = 0 = (cid:118)∫ J i da + C 0→ C S d dt S (cid:118)∫ ε 0E i da dVρ ∫ V 6.641, Electromagnetic Fields, Forces, and Motion Prof. Markus Zahn Lecture 1 Page 2 of 6 J i da + (cid:118)∫ S d dt ∫ V ρ dV = 0 Total current leaving volume inside volume through surface Total charge 6. Lorentz Force Law f = q E + v × µ H) ( 0 II. Electric Field from Point Charge (cid:118)∫ S ε0E i da = ε0E 4 π r = q 2 r E = r q 4π ε0r 2 T sin θ = f = c 2 q 4π ε0r 2 T cos θ = Mg tan θ = 2 q 4π ε0r 2 Mg = r 2l 6.641, Electromagnetic Fields, Forces, and Motion Prof. Markus Z
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r 2 Mg = r 2l 6.641, Electromagnetic Fields, Forces, and Motion Prof. Markus Zahn Lecture 1 Page 3 of 6 3 Mg ⎡2π ε0r q = ⎢ l ⎣ 2 1 ⎤ ⎥ ⎦ III. Faraday Cage J da = i = - i ∫(cid:118) S d dt ∫ ρ dV = - d dt ( ) = -q dq dt ∫ idt = q 6.641, Electromagnetic Fields, Forces, and Motion Prof. Markus Zahn Lecture 1 Page 4 of 6 IV. Boundary Conditions 1. Gauss’ Continuity Condition Courtesy of Krieger Publishing. Used with permission. ε0E i da = ∫σsdS ⇒ ε0 (E2n (cid:118)∫ - E1n ) dS = σ dS s S S ε0 (E - E ) = σ ⇒ n i ⎡ε0 (E - E1 )⎤ = σ 1n 2n 2 s s ⎦ ⎣ 2. Continuity of Tangential E Courtesy of Krieger Publishing. Used with permission. E i ds = (E - E2t ) dl = 0 ⇒ E - E2t = 0 1t 1t (cid:118)∫ C n× E1 - E2 ) = 0 ( Equivalent to = Φ2 along boundary Φ1 6.641, Electromagnetic Fields, Forces, and Motion Prof. Markus Zahn Lecture 1 Page 5 of 6 3. Normal H ∇ µ H = 0 ⇒ i 0 (cid:118)∫ µ0 S H i da = 0 4. Tangential H µ (H - Hbn an 0 ) A = 0 H = Hbn
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Tangential H µ (H - Hbn an 0 ) A = 0 H = Hbn an n i ⎡ ⎣ H a - H b ⎤ = 0 ⎦ (cid:118) ∇ × H = J ⇒ ∫ H i ds = J i da ∫ C S H ds - Hatds = Kds bt H - Hat = K bt n × ⎡ ⎣ H a - H b ⎤ = K ⎦ ∂ρ ∇ i J + = 0 ∂t 5. Conservation of Charge Boundary Condition i(cid:118)∫ d ∫ J da + ρdV = 0 dt V S n i ⎣ ⎡ J a - Jb ⎦ ⎤ + ∂ t ∂ σ = 0 s 6.641, Electromagnetic Fields, Forces, and Motion Prof. Markus Zahn Lecture 1 Page 6 of 6
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Kepler’s Second Law By studying the Danish astronomer Tycho Brahe’s data about the motion of the planets, Kepler formulated three empirical laws; two of them can be stated as follows: Second Law A planet moves in a plane, and the radius vector (from the sun to the planet) sweeps out equal areas in equal times. First Law The planet’s orbit in that plane is an ellipse, with the sun at one focus. From these laws, Newton deduced that the force keeping the planets in their orbits had magnitude 1/d2, where d was the distance to the sun; moreover, it was directed toward the sun, or as was said, central, since the sun was placed at the origin. Using a little vector analysis (without coordinates), this section is devoted to showing that the Second Law is equivalent to the force being central. It is harder to show that an elliptical orbit implies the magnitude of the force is of the form K/d2, and vice-versa; this uses vector analysis in polar coordinates and requires the solution of non-linear differential equations. 1. Differentiation of products of vectors Let r(t) and s(t) be two differentiable vector functions in 2- or 3-space. Then (1) d dt (r · s) = dr dt · s + r · ds ; dt d dt (r × s) = dr dt × s + r ×
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ds ; dt d dt (r × s) = dr dt × s + r × ds . dt These rules are just like the product rule for differentiation. Be careful in the second 6 b × a in general. rule to get the multiplication order correct on the right, since a × b = The two rules can be proved by writing everything out in terms of i , j , k components and differentiating. They can also be proved directly from the definition of derivative, without resorting to components, as follows: Let t increase by Δt. Then r increases by Δr, and s by Δs, and the corresponding change in r · s is given by Δ(r · s) = (r + Δr) · (s + Δs) − r · s , so if we expand the right side out and divide all terms by Δt, we get Δ(r · s) Δt = Δr Δt · s + r · Δs Δt + Δr Δt · Δs . Now let Δt → 0; then Δs → 0 since s(t) is continuous, and we get the first equation in (1). The second equation in (1) is proved the same way, replacing · by × everywhere. 2. Kepler’s second law and the central force. To show that the force being central (i.e., directed toward the sun) is equivalent to Kepler’s second law,
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being central (i.e., directed toward the sun) is equivalent to Kepler’s second law, we need to translate that law into calculus. “Sweeps out equal areas in equal times” means: the radius vector sweeps out area at a constant rate . 1 2 KEPLER’S SECOND LAW The first thing therefore is to obtain a mathematical expression for this rate. Referring to the picture, we see that as the time increases from t to t + Δt, the corresponding change in the area A is given approximately by ΔA ≈ area of the triangle = 1 2 |r × Δr| , since the triangle has half the area of the parallelogram formed by r and Δr; thus, 2 ΔA Δt � � � Δr � � � � ≈ r × , � Δt and as Δt → 0, this becomes (2) 2 dA dt � � � = r × � � dr � � = � dt |r × v|. where v = dr dt . r + Δ r Δ r r Using (2), we can interpret Kepler’s second law mathematically. Since the area is swept out at a constant rate, dA/dt is constant, so according to (2), (3) |r × v| is a constant. Moreover, since Kepler’s law says r lies in a plane, the velocity vector v also lies in the same plane, and therefore (4) r × v has constant direction (perpendicular to the plane of motion).
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same plane, and therefore (4) r × v has constant direction (perpendicular to the plane of motion). Since the direction and magnitude of r × v are both constant, (5) r × v = K, a constant vector, and from this we see that (6) d dt (r × v) = 0 . But according to the rule (1) for differentiating a vector product, (7) d dt (r × v) = v × v + r × a, where a = dv , dt = r × a, since s × s = 0 for any vector s. Now (6) and (7) together imply (8) r × a = 0, which shows that the acceleration vector a is parallel to r, but in the opposite direction, since the planets do go around the sun, not shoot off to infinity. Thus a is directed toward the center (i.e., the sun), and since F = ma, the force F is also directed toward the sun. (Note that “center” does not mean the center of the elliptical orbits, but the mathematical origin, i.e., the tail of the radius vector r, which we are taking to be the sun’s position.) The reasoning is reversible, so for motion under any type of central force, the path of motion will lie in a plane and area will be swept out by the radius vector at a constant rate. MIT Open
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in a plane and area will be swept out by the radius vector at a constant rate. MIT OpenCourseWare http://ocw.mit.edu 18.02SC Multivariable Calculus Fall 2010 For information about citing these materials or our Terms of Use, visit: http://ocw.mit.edu/terms.
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55 1.25. The Quantum Group Uq (sl2). Let us consider the Lie algebra sl2. Recall that there is a basis h, e, f ∈ sl2 such that [h, e] = 2e, [h, f] = −2f, [e, f] = h. This motivates the following definition. Definition 1.25.1. Let q ∈ k, q =� ±1. The quantum group U (sl2) is generated by elements E, F and an invertible element K with defining relations q KEK −1 = q 2E, KFK −1 = q−2F, [E, F] = K − K −1 . q − q−1 Theorem 1.25.2. There exists a unique Hopf algebra structure on Uq (sl2), given by • Δ(K) = K ⊗ K • Δ(E) = E ⊗ K + 1 ⊗ E; • Δ(F) = F ⊗ 1 + K −1 ⊗ F (thus K is a grouplike element); (thus E, F are skew-primitive ele­ ments). Exercise 1.25.3. Prove Theorem 1.25.2. Remark 1.25.4. Heuristically, K = qh, and thus K − K −1 lim 1→ q − q−1 q = h. → 1, the relations of Uq (sl2) degenerate into the So in the limit q relations of U (sl2), and thus Uq (sl2) should be viewed as a Hopf algebra deformation of the enveloping algebra U (sl2). In fact, one can make this heuristic idea into a precise statement, see e.g. [K]. If q is a root of unity, one can also define a finite dimensional version of Uq (sl2).
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K]. If q is a root of unity, one can also define a finite dimensional version of Uq (sl2). Namely, assume that the order of q is an odd number �. Let uq (sl2) be the quotient of Uq (sl2) by the additional relations E� = F� = K � − 1 = 0. Then it is easy to show that uq (sl2) is a Hopf algebra (with the co­ product inherited from Uq (sl2)). This Hopf algebra is called the small quantum group attached to sl2. 1.26. The quantum group Uq (g). The example of the previous sub­ section can be generalized to the case of any simple Lie algebra. Namely, let g be a simple Lie algebra of rank r, and let A = (aij ) be its Cartan matrix. Recall that there exist unique relatively prime positive integers di, i = 1, . . . r such that diaij = dj aji. Let q ∈ k, q =� ±1. Definition 1.26.1. • The q-analog of n is nq − q−n . q − q−1 [n]q = 56 • The q-analog of the factorial is [n]q ! = n � [l]q = l=1 (q − q−1) · · · (qn − q−n) . (q − q−1)n Definition 1.26.2. The quantum group Uq(g) is generated by elements Ei, Fi and invertible elements Ki, with defining relations KiKj = Kj Ki, KiEj Ki −1 = q aij Ej , KiFj Ki −1 = q−aij Fj , [Ei, Fj ] = δij K di − K −di i qdi − q−di i , and the q-Serre relations: (1.26.1) and (1.26
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− q−di i , and the q-Serre relations: (1.26.1) and (1.26.2) 1−aij � [l]q ![1 i l=0 l 1) − ( a − ij 1−aij −lEj Ei Ei l = 0, i =� − l]q ! i 1−aij � [l]q ![1 i l=0 l 1) − ( a − ij 1−aij −lFj Fl Fi i = 0, i =� − l]q ! i j j. More generally, the same definition can be made for any symmetriz­ able Kac-Moody algebra g. Theorem 1.26.3. (see e.g. [CP]) There exists a unique Hopf algebra structure on Uq (g), given by • Δ(Ki) = Ki ⊗ Ki; • Δ(Ei) = Ei ⊗ Ki + 1 ⊗ Ei; • Δ(Fi) = Fi ⊗ 1 + Ki −1 ⊗ Fi. Remark 1.26.4. Similarly to the case of sl2, in the limit q 1, these relations degenerate into the relations for U (g), so Uq(g) should be viewed as a Hopf algebra deformation of the enveloping algebra U (g). → 1.27. Categorical meaning of skew-primitive elements. We have seen that many interesting Hopf algebras contain nontrivial skew-primitive elements. In fact, the notion of a skew-primitive element has a cate­ gorical meaning. Namely, we have the following proposition. Proposition 1.27.1. Let g, h be grouplike elements of a coalgebra C, and Primh,g(C) be the space of skew-primitive elements of type
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coalgebra C, and Primh,g(C) be the space of skew-primitive elements of type h, g. Then the space Primh,g(H)/k(h − g) is naturally isomorphic to Ext1(g, h), where g, h are regarded as 1-dimensional right C-comodules. Proof. Let V be a 2-dimensional H-comodule, such that we have an exact sequence 0 → → → → V g h 0. 57 Then V has a basis v0, v1 such that π(v0) = v0 ⊗ h, π(v1) = v1 ⊗ x + v0 ⊗ g. The condition that this is a comodule yields that x is a skew-primitive element of type (h, g). So any extension defines a skew-primitive el­ v0,→ ement and vice versa. Also, we can change the basis by v0 v1 + λv0, which modifies x by adding a trivial skew-primitive v1 � element. This implies the result. → Example 1.27.2. The category C of finite dimensional comodules over uq (sl2) is an example of a finite tensor category in which there are objects V such that V ∗∗ is not isomorphic to V . Namely, in this category, the functor V �→ V ∗∗ is defined by the squared antipode S2 , which is conjugation by K: S2(x) = KxK −1 . Now, we have Ext1(K, 1) = Y = �E, F K�, a 2-dimensional space. The set of iso­ morphism
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Y = �E, F K�, a 2-dimensional space. The set of iso­ morphism classes of nontrivial extensions of K by 1 is therefore the projective line PY . The operator of conjugation by K acts on Y with eigenvalues q2, q−2, hence nontrivially on PY . Thus for a generic ex­ tension V , the object V ∗∗ is not isomorphic to V . However, note that some power of the functor ∗∗ on C is isomorphic (in fact, monoidally) to the identity functor (namely, this power is the order of q). We will later show that this property holds in any finite tensor category. Note also that in the category C, V ∗∗ ∼= V if V is simple. This clearly has to be the case in any tensor category where all simple objects are invertible. We will also show (see Proposition 1.41.1 below) that this is the case in any semisimple tensor category. An example of a tensor category in which V ∗∗ is not always isomorphic to V even for simple V is the category of finite dimensional representations of the the Yangian H = Y (g) of a simple complex Lie algebra g, see [CP, 12.1]. Namely, for any finite dimensional representation V of H and any complex number z one can define the shifted representation V (z) (such that V (0) = V ). Then V ∗∗ ∼ V (2h∨), where h∨ is the dual Coxeter number of g,
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�∗ ∼ V (2h∨), where h∨ is the dual Coxeter number of g, see [CP, p.384]. If V is a non-trivial irreducible finite dimensional representation then V (z) = V for z = 0. Thus, V ∗∗ ∼� = V . Moreover, we see that the functor ∗∗ has infinite order even when restricted to simple objects of C. = ∼ However, the representation category of the Yangian is infinite, and the answer to the following question is unknown to us. Question 1.27.3. Does there exist a finite tensor category, in which there is a simple object V such that V ∗∗ is not isomorphic to V ? (The answer is unknown to the authors). � � 58 Theorem 1.27.4. Assume that k has characteristic 0. Let C be a finite ring category over k with simple object 1. Then Ext1(1, 1) = 0. Proof. Assume the contrary, and suppose that V is a nontrivial exten­ sion of 1 by itself. Let P be the projective cover of 1. Then Hom(P, V ) is a 2-dimensional space, with a filtration induced by the filtration on V , and both quotients naturally isomorphic to E := Hom(P, 1). Let v0, v1 be a basis of Hom(P, V ) compatible to the filtration, i.e. v0 spans the 1-dimensional subspace defined by the filtration. Let A = End(P ) (this is a finite dimensional algebra). Let ε : A k be the character defined by the (right) action of A on E. Then the matrix of a ∈ A in the basis v0, v1 has the form →
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(right) action of A on E. Then the matrix of a ∈ A in the basis v0, v1 has the form → [a]1 = � ε(a) χ1(a) ε(a) � (1.27.1) 0 where χ1 ∈ A∗ is nonzero. Since a → [a]1 is a homomorphism, χ1 is a derivation: χ1(xy) = χ1(x)ε(y) + ε(x)χ1(y). Now consider the representation V ⊗ V . Using the exactness of tensor products, we see that the space Hom(P, V ⊗V ) is 4-dimensional, and has a 3-step filtration, with successive quotients E, E ⊕ E, E, and basis v00; v01, v10; v11, consistent with this filtration. The matrix of a ∈ End(P ) in this basis is (1.27.2) ⎛ [a]2 = ⎜ ⎜ ⎝ ⎞ ε(a) χ1(a) χ1(a) χ2(a) χ1(a) ⎟ ⎟ ⎠ ε(a) χ1(a) ε(a) ε(a) 0 0 0 0 0 0 0 Since a → [a]2 is a homomorphism, we find χ2(ab) = ε(a)χ2(b) + χ2(a)ε(b) + 2χ1(a)χ1(b). We can now proceed further (i.e. consider V ⊗V ⊗V etc.) and define for every positive n, a linear function χn ∈ A∗ which satisfies the equation � � � n j χj (a)χn−j (b), χn(ab) = n j=0 ε. � χm(a)s Thus for any s ∈ k, we can
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