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Customer Service - Understanding your water meter If you have an AMI meter click this link to the Customer Service Portal (My Account) to see your usage analytics. Otherwise information on your analog meter follows. Understanding and Reading Your Analog Meter There are two basic types of analog water meters – straight reading and circular reading. The straight reading meter records cubic feet of water used in much the same way that a car’s odometer records miles. The dial with a single hand measures tenth of a cubic foot. The circular reading meter uses a series of circular dials to record cubic feet of water used. To read this type of meter, start with the 100,000th circle, then read the 10,000 circle, and so forth on down to the circle that reads in 1 cubic foot increments. If a hand is between two numbers, always read the lower number. This meter, for example, shows a reading of 2,425.92 cubic feet (the “6” in the last dial position has not quite rolled over). One cubic foot is 7.48 gallons of water. This meter has registered almost 18,146 gallons since it was new (2,425.92 x 7.48). When the Water Authority reads your meter, only the white dials with the black letters are read to measure the number of “units” that have been used. One unit equals 100 cubic feet, or 748 gallons. In this example, the meter reading would be 24. This reading appears on the left side of your bill, in the box labeled “Metered Usage.” If the meter reading the month prior was 20, this customer would be billed for four units. The triangle is a leak detector. The triangle will rotate if water is passing through the meter. If no one is using water, but the triangle is turning, you may have an undiscovered leak in your plumbing system. For more information on your water usage: Understanding and Measuring Water Usage & Checking for Leaks Understanding and Measuring Water Usage To measure the amount of water used for any activity, follow these instructions: 1. Before you begin the measurement, write down the meter reading to two decimal places. (see the photo example above for help). 2. Perform the activity you want to measure, but be sure that no other water is being used during the test. Some examples include: • Washing a load of laundry or dishes • Taking a shower or bath • Washing your car • Watering your lawn • Filling your swimming pool 3. If you are doing an activity that may have a variable duration, such as taking a shower or running your sprinklers, you should measure the number of minutes the activity required. This information will allow you to determine the number of gallons per minute the activity requires. 4. After the activity is complete, read the meter again. Subtract the first reading from the second reading, and then multiply the remainder by 7.48 to convert to gallons. 5. To get gallons per minute, divide the number of gallons by the number of minutes the activity required. For example, if you ran your sprinkler system for 10 minutes and used 110 gallons of water, your system uses 11 gallons per minute! Keep in mind that the “units” shown on your water bill are equal to 100 cubic feet. The Water Authority water meter readers only record complete units for billing purposes, so the last two digits (the “tens” and the “ones” figures) are omitted. Remember: one unit equals 748 gallons, so one cubic foot equals about 7.5 gallons. Using your Meter to Check for Leaks If you suspect you have a leak, you can measure the volume: • Write down the meter reading and the time of day to the minute. • Don’t use any water during the test. Usually it is best to do this when you will be away from home for an hour or more. Make sure devices such as evaporative coolers and ice makers are turned • Read the meter again when you return and note the time of day. • Subtract the second reading from the first. Multiply the remainder by 7.48. This is the number of gallons that passed through the meter during the test period. • Divide the amount of water by the number of minutes in the test. For example, if 17 gallons leaked out during a 180 minute period, you have a leak of 0.094 gallons per minute. • Multiply the gallons per minute by 1,440 to calculate gallons per day. Multiply gallons per minute by 43,920 to calculate gallons per month. In this example, just 0.094 gallons per minute equates to over 4,128 gallons each month!
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6 GMAT Function Practice Questions With Explanations // Ambitio 2 September 2024 6 minutes read 6 GMAT Function Practice Questions With Explanations Key Takeaways • Focus on understanding functions, their compositions, and how to solve them. • Regular practice with real-world examples is crucial for mastering GMAT function questions. • Employ strategy like visualization, and the two-pass system to enhance your learnings. • Teach the material to others and engage actively to solidify your knowledge. • Practice under real test conditions to build stamina and familiarity with the test environment. Preparing for the GMAT test? We understand your headache. GMAT exams can be challenging, especially when it comes to mastering functions and equations. In this last-minute guide, we will cover 6 GMAT function practice questions that will enhance your preparation. Understanding the value of each variable, how to plug numbers into equations, and defining the domain are crucial skills for excelling in the GMAT. We provide detailed explanations and expert solutions to ensure you grasp each concept thoroughly. Whether you’re working with integers, sequences, or complex functions, these practice questions are designed to help you succeed. So, stick to the guide till the end – we gonna cover the essential elements of GMAT functions to boost your score and confidence. Worried about the cost of Studying Abroad? Sign up to access 25 game-changing scholarships that could cover your costs. What are the function questions in the GMAT exam? GMAT function questions are a key part of the quantitative section that requires a solid understanding of mathematical expressions and their applications. These questions often involve interpreting functions, using parentheses correctly, and solving for the exact value of variables like “x”. Practice tests are invaluable for mastering these concepts because they provide real-world examples and challenges similar to those found on the actual GMAT. In addition to the GMAT test prep, getting ready for other standardized tests like the SAT, ACT test prep, and SSAT also benefits from a strong grasp of functions. Whether you’re enrolled in ISEE courses, SSAT test prep, or even MCAT courses, understanding functions is crucial. For those seeking personalized guidance, options abound from tutors to classes in various locations. By honing your skills through targeted practice, you can achieve precise output and improve your overall performance on test day. Stuck on How to Pick Your Ideal College? Sign up to access your tailored shortlist and simplify finding your ideal college. 6 mostly asked GMAT function practice questions with explanations The following questions cover various concepts related to functions, including function composition, domain and range, inverse functions, and properties of exponential and logarithmic functions. Here are 6 mostly-asked GMAT function practice questions related to functions: 1. If f(x) = 2x + 3 and g(x) = x^2 – 2, what is f(g(2))? Evaluate π (2)g(2): The function π (π ₯)g(x) is defined as: π (π ₯)=π ₯2β 2g(x)=x2β 2 We substitute π ₯=2x=2 into the function π (π ₯)g(x): π (2)=22β 2g(2)=22β 2 Calculate the exponent and subtraction: π (2)=4β 2=2g(2)=4β 2=2 Evaluate π (π (2))f(g(2)): We now have π (2)=2g(2)=2. Next, we use this result as the input for the function π (π ₯)f(x). The function π (π ₯)f(x) is defined as: π (π ₯)=2π ₯+3f(x)=2x+3 We substitute π ₯=2x=2 (which is the result of π (2)g(2)) into the function π (π ₯)f(x): π (2)=2(2)+3f(2)=2(2)+3 Perform the multiplication and addition: π (2)=4+3=7f(2)=4+3=7 Therefore, π (π (2))=7f(g(2))=7. 2. Let h(x) = (x^2 + 3x – 2) / (x – 1). Find the value of h(2). To find the value of β (2)h(2) for the function β (π ₯)=π ₯2+3π ₯β 2π ₯β 1h(x)=xβ 1×2+3xβ 2β , we need to substitute π ₯=2x=2 into the function and simplify. • Substitute π ₯=2x=2 into β (π ₯)h(x): β (2)=22+3(2)β 22β 1h(2)=2β 122+3(2)β 2β • Simplify the numerator: 22=422=4 3(2)=63(2)=6 4+6β 2=84+6β 2=8 • Simplify the denominator: 2β 1=12β 1=1 • Combine the simplified numerator and denominator: β (2)=81=8h(2)=18β =8 Therefore, the value of β (2)h(2) is 8. 3. If f(x) = 3x^2 – 2x + 5 and g(x) = 2x – 1, find (f β g)(x). To find (π β π )(π ₯)(fβ g)(x), which is the composition of the functions π (π ₯)f(x) and π (π ₯)g(x), we need to substitute π (π ₯)g(x) into π (π ₯)f(x). This means we will replace every π ₯x in π (π ₯)f(x) with π (π ₯)g(x). Given: π (π ₯)=3π ₯2β 2π ₯+5f(x)=3×2β 2x+5 π (π ₯)=2π ₯β 1g(x)=2xβ 1 We want to find π (π (π ₯))f(g(x)). • Substitute π (π ₯)g(x) into π (π ₯)f(x): π (π (π ₯))=π (2π ₯β 1)f(g(x))=f(2xβ 1) • Replace every π ₯x in π (π ₯)f(x) with 2π ₯β 12xβ 1: π (2π ₯β 1)=3(2π ₯β 1)2β 2(2π ₯β 1)+5f(2xβ 1)=3(2xβ 1)2β 2(2xβ 1)+5 • Expand and simplify (2π ₯β 1)2(2xβ 1)2: (2π ₯β 1)2=(2π ₯β 1)(2π ₯β 1)=4π ₯2β 4π ₯+1(2xβ 1)2=(2xβ 1)(2xβ 1)=4×2β 4x+1 • Substitute back into the function: π (2π ₯β 1)=3(4π ₯2β 4π ₯+1)β 2(2π ₯β 1)+5f(2xβ 1)=3(4×2β 4x+1)β 2(2xβ 1)+5 • Distribute the constants: π (2π ₯β 1)=12π ₯2β 12π ₯+3β 4π ₯+2+5f(2xβ 1)=12×2β 12x+3β 4x+2+5 • Combine like terms: π (2π ₯β 1)=12π ₯2β 16π ₯+10f(2xβ 1)=12×2β 16x+10 Therefore, (π β π )(π ₯)=12π ₯2β 16π ₯+10(fβ g)(x)=12×2β 16x+10. 4. The function f is defined by f(x) = 2^x for all real numbers x. Find f(log2 8). To find π (logβ ‘28)f(log2β 8) for the function π (π ₯)=2π ₯f(x)=2x, follow these steps: • Evaluate the inner expression logβ ‘28log2β 8: The logarithm logβ ‘28log2β 8 asks the question: “To what power must 2 be raised to get 8?” Since 23=823=8: logβ ‘28=3log2β 8=3 • Substitute the value of logβ ‘28log2β 8 into π (π ₯)f(x): Given π (π ₯)=2π ₯f(x)=2x, we need to find π (3)f(3): π (3)=23f(3)=23 • Calculate 2323: 23=823=8 Therefore, π (logβ ‘28)=8f(log2β 8)=8. 5. Let f(x) = |x – 3| and g(x) = 2x – 1. Find the value(s) of x for which f(x) = g(x). To find the value(s) of π ₯x for which π (π ₯)=π (π ₯)f(x)=g(x), given π (π ₯)=β £π ₯β 3β £f(x)=β £xβ 3β £ and π (π ₯)=2π ₯β 1g(x)=2xβ 1, we need to solve the equation β £π ₯β 3β £=2π ₯β 1β £xβ 3β £=2xβ 1. The absolute value equation β £π ₯β 3β £=2π ₯β 1β £xβ 3β £=2xβ 1 can be split into two separate equations: 1. Case 1: π ₯β 3=2π ₯β 1xβ 3=2xβ 1 2. Case 2: β (π ₯β 3)=2π ₯β 1β (xβ 3)=2xβ 1 Let’s solve each case separately. Case 1: π ₯β 3=2π ₯β 1xβ 3=2xβ 1 • Subtract π ₯x from both sides: β 3=π ₯β 1β 3=xβ 1 • Add 1 to both sides: β 2=π ₯β 2=x So, π ₯=β 2x=β 2.Case 2: β (π ₯β 3)=2π ₯β 1β (xβ 3)=2xβ 1 • Distribute the negative sign: β π ₯+3=2π ₯β 1β x+3=2xβ 1 • Add π ₯x to both sides: 3=3π ₯β 13=3xβ 1 • Add 1 to both sides: 4=3π ₯4=3x • Divide by 3: π ₯=43x=34β • Check π ₯=β 2x=β 2: π (β 2)=β £β 2β 3β £=β £β 5β £=5f(β 2)=β £β 2β 3β £=β £β 5β £=5 π (β 2)=2(β 2)β 1=β 4β 1=β 5g(β 2)=2(β 2)β 1=β 4β 1=β 5 Since π (β 2)β π (β 2)f(β 2)ξ =g(β 2), π ₯=β 2x=β 2 is not a solution. • Check π ₯=43x=34β : π (43)=β £43β 3β £=β £43β 93β £=β £β 53β £=53f(34β )=β £β £β 34β β 3β £β £β =β £β £β 34β β 39β β £β £β =β £β £β β 35β β £β £β =35β π (43)=2(43)β 1= 83β 33=53g(34β )=2(34β )β 1=38β β 33β =35β Since π (43)=π (43)f(34β )=g(34β ), π ₯=43x=34β is a solution. Therefore, the value of π ₯x for which π (π ₯)=π (π ₯)f(x)=g(x) is 4334β . 6. If f(x) = 3^(2x) and g(x) = log3(x^2 – 1), find f(g(8)). To find π (π (8))f(g(8)) for the functions π (π ₯)=32π ₯f(x)=32x and π (π ₯)=logβ ‘3(π ₯2β 1)g(x)=log3β (x2β 1), we need to follow these steps: • Evaluate π (8)g(8): The function π (π ₯)g(x) is defined as: π (π ₯)=logβ ‘3(π ₯2β 1)g(x)=log3β (x2β 1) Substitute π ₯=8x=8 into π (π ₯)g(x): π (8)=logβ ‘3(82β 1)g(8)=log3β (82β 1) Calculate the value inside the logarithm: 82=6482=64 64β 1=6364β 1=63 So, π (8)=logβ ‘3(63)g(8)=log3β (63) • Evaluate π (π (8))f(g(8)): Now we need to use this result as the input for the function π (π ₯)f(x). The function π (π ₯)f(x) is defined as: π (π ₯)=32π ₯f(x)=32x Substitute π ₯= logβ ‘3(63)x=log3β (63) into π (π ₯)f(x): π (logβ ‘3(63))=32β logβ ‘3(63)f(log3β (63))=32β log3β (63) • Simplify the expression: Using the property of logarithms and exponents: 32β logβ ‘3(63)=3logβ ‘3(632)32β log3β (63)=3log3β (632) Since 3logβ ‘3(π )=π 3log3β (a)=a: 3logβ ‘3(632)= 6323log3β (632)=632 • Calculate 632632: 632=3969632=3969 Therefore, π (π (8))=3969f(g(8))=3969. See how Successful Applications Look Like! Access 350K+ profiles of students who got in. See what you can improve in your own application! Tips to answer better in GMAT exams Look – we know that the GMAT is hard – especially if you are from a non-math background, but winning test-taking strategies can make a significant difference in your performance. Beyond the standard advice of regular practice and time management, unique and insightful tips can give you an edge. Here are seven uncommon strategies to help you excel on exam day. Visualize Success Before diving into your study session or the exam itself, take a few minutes to close your eyes and visualize yourself successfully answering questions. This mental rehearsal can boost your confidence and reduce anxiety, making you more focused and effective during the test. Use the Elimination Method Instead of looking for the correct answer right away, start by eliminating the obviously wrong choices. This strategy can increase your chances of selecting the right answer by narrowing down your options, especially in tricky quantitative and verbal questions. Practice Mindfulness and Breathing Techniques Incorporate mindfulness exercises and deep breathing into your study routine. These techniques can help you stay calm and maintain concentration during the exam, particularly during challenging sections or when facing time pressure. Teach the Material Another way of learning and enhancing our knowledge retention is through facilitated learning where one is encouraged to pass the knowledge being studied. Student-Directed Activities: Look for a study buddy or practice explaining what you learned to an imaginary audience. This approach helps to consolidate information within you and show a focus on aspects that require more focus. Use the Two-Pass System During the exam, go through the entire section quickly first, answering the questions you find easiest. On the second pass, tackle the more difficult questions. This approach ensures you secure easy points and manage your time more effectively. Develop a Question Identification System Create a personal system to quickly identify the type of question you’re facing (e.g., algebra, geometry, critical reasoning). This system allows you to switch mental gears efficiently and apply the most appropriate strategies for each question type. Simulate Test Conditions Regularly practice under actual test conditions. Use a timer, sit in a quiet environment, and take full-length practice tests. Simulating the test environment helps you build stamina and get accustomed to the pressure and timing constraints of the GMAT. Start Your University Applications with Ambitio Pro! Get Ambitio Pro! Begin your journey to top universities with Ambitio Pro. Our premium platform offers you tools and support needed to craft standout applications. Unlock Advanced Features for a More Comprehensive Application Experience! Start your Journey today Start your Journey today Remember to practice regularly, use strategic approaches, and stay calm under pressure. These practice questions and tips are designed to help you excel and achieve your desired GMAT score. To enhance your understanding, actively engage with the material. Instead of passively reading, try solving problems without looking at the solutions first. Discuss questions with peers, and donβ t hesitate to teach others. Active learning solidifies your knowledge and exposes areas needing improvement, ultimately leading to better performance on test day. Transform your GMAT preparation with Ambitio’s expert guidance. Our comprehensive approach includes personalized study plans, adaptive practice tests, and strategic insights, all designed to enhance your understanding and performance across the exam’s quantitative and verbal sections. Stuck on How to Pick Your Ideal College? Sign up to access your tailored shortlist and simplify finding your ideal college. How do I Register for the GMAT? How Much Does the GMAT Cost? How Often Can I Take the GMAT? You can take the GMAT up to five times every 12 months, with a limit of not more than once in a 16-day period or more than eight times in total Can I Reschedule a GMAT Appointment? You can reschedule your GMAT appointment by logging into your personal GMAT account at mba.com, but there are fees involved depending on the timing of the rescheduling Can I Retake the GMAT? What Material Is Tested on the GMAT? The GMAT tests essential skills needed in business school and subsequent business careers. It includes sections like analytical writing assessment, integrated reasoning, quantitative, and verbal, each assessing different skills and concepts Table of Contents Almost there! Just enter your OTP, and your planner will be on its way! Code sent on Resend OTP (30s) Your Handbook Is Waiting on WhatsApp! 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Teaching the Types of Triangles with Math Wheel Magic! - Cognitive Cardio MathCognitive Cardio Math Have you noticed how triangles can be a bit of a math maze for our students? I mean, equilateral, isosceles, scalene – it’s like they’ve entered a geometric wonderland without a compass! I have felt this way time and again, so I set out to create a tool to help break down the different types of triangles in a memorable way for my students. I’m excited to share with you my student-approved Types of Triangles Math Wheel Guided Notes! It’s the treasure map your students need to navigate the triangle jungle. With color-coded notes, practice problems scattered around the wheel, and a graphic organizer that breaks down and organizes the different types, you’ll be turning confusion into a full-blown math adventure. Get ready for those “Aha!” moments that turn the triangle saga from puzzling to downright fun! What Are Math Doodle Wheels? I’m standing at the front of my class and all I see are blank stares, math nerves, and a whole bunch of “I’d-rather-be-anywhere-else” vibes. My students were going through the motions of note-taking, but it was clear, the info wasn’t sticking. That’s when the lightbulb moment happened, and Math Doodle Wheels were born! I wanted notes to be more than just a checkmark on the to-do list. I wanted them to be a trusty sidekick, a go-to guide during units and review sessions. So, I crafted these wheels where math concepts are broken down into bite-size sections of the wheel, vocab gets its VIP definition moment, and each step is practiced with examples. Plus, there’s space for some extra practice. The secret? We’re color-coding and doodling symbols as we go to help them remember the steps and vocabulary. Suddenly, note-taking became an art, and my kiddos enjoyed it and used these notes as tools. To learn more about Doodle Wheels and their versatility, check out Transform Your Upper Elementary or Middle School Math Class with this Unique Note-Taking Method! How I Teach Types of Triangles We’re diving headfirst into the world of triangles, and it all begins right at the heart of our Types of Triangles Math Wheel. I gather my students around, and we start chatting about the six main players in the triangle game – equilateral, scalene, acute, obtuse, right, and isosceles. Now, here’s the math magic. No matter how different these triangles are, their angles always add up to 180°. I have my students write that 180° right in the center of our wheel to be a constant reminder. With the center complete, I kick off the grand tour of the wheel. The first thing we discuss is that there are two ways we can classify triangles. We can classify them by their sides and we can classify them by their angles. To keep this in mind, we tackle the three types of triangles classified by their sides first, then we move on to the three types of triangles classified by angles. With that explanation, we have now taken 6 types of triangles and grouped them into two groups of three and made them a little easier to understand. Sections of the Types of Triangles Equilateral Triangle I explain that this triangle is known for having three equal sides, which we jot down in that section. I draw three hash marks, one on each leg of the triangle, to show that they are equal. To give them an example of number measurements, I write 6 meters on all three sides. I like starting with equilateral because it’s the least intimidating since all three sides are the same Scalene Triangle Now, onto the scalene triangle. I kind of think of it as the rebel of the triangle crew. It’s the complete opposite of our equilateral triangle because it plays by its own rules with no equal sides. We fill in the blank in our notes section. Then, it’s example time. I draw my trusty triangle and label it with the following measurements: 5 cm on the short leg, 7 cm on the top leg, and 10 cm on the bottom leg. Why? To show my students how each side has its unique length. Isosceles Triangle Now, we wrap up our note-taking adventure with the isosceles triangle. An isosceles triangle is all about the buddy system, boasting two equal sides. In our triangle example, I draw single hash marks on the left and right legs because they’re equal measurements. I label both legs with 4 feet, emphasizing their equal status. Now, I draw two hash marks on the bottom leg, signaling that its measurements are playing solo. I label 6 feet for the bottom leg, driving home the point that it’s not on the equal side with the left and right legs. Acute Triangle Next up is the acute triangle. I walk my students through jotting down the key detail that all three angles inside an acute triangle are less than 90°. I draw two examples to show my students a couple of different scenarios they may encounter. In the first triangle, I sketch a curved line at each corner, representing an angle, and drop 60° in each. We add each of the three angles, each at 60°, sum up to 180°, sealing the deal that triangles have a total of 180° inside. I draw another triangle to explain that acute triangles won’t always have the same angle measurements. In this second one, I draw those curved lines again and jot down 70°, 60°, and 50° in the corners. Quick adding by my eager students confirms the angles total to 180°. Obtuse Triangle In our next section, we dive into the world of the obtuse triangle. For a triangle to earn the “obtuse” badge, it needs to boast one obtuse angle, meaning one angle must be greater than 90°. In my example, I draw an obtuse triangle. I draw the curved line to represent the angle in the top angle inside the triangle and write 120°. We talk about how we can check that this is obtuse by noting how it’s a number larger than 90. Right Triangle We move to the right triangle section. We review what makes a right angle, which is when a 90° angle is formed at the intersection of two perpendicular lines. I have my students write down that one angle must measure 90°. Then, in the triangle, we practice drawing the square symbol where the two lines intersect to represent 90°. What’s Next? Once we wrap up our note-taking session, it’s time to put our knowledge to the test by tackling the examples scattered around the math doodle wheel. I pre-select a few examples, modeling to my students my thought process and reasonings for classifying which triangles to which measurements. After a couple of examples, I give the steering wheel to my students. They become the teacher and start leading me through the steps for a few examples. Then, I partner them up to complete the remaining practice examples. I set a timer so we have time to check their answers and talk through their reasoning. After conquering those practice problems, I give my math enthusiasts a little artistic freedom. They go back into their notes, adding splashes of color and extra symbols to etch the differences between the triangles into their memory. It’s a personal touch that not only triggers those “aha” moments but also makes math feel like a comfortable space. The math doodle wheel isn’t about bombarding them with gibberish and 50 problems. It’s all about breaking the skill into bite-sized chunks, creating a handy reference tool that’ll be their math companion throughout the year! More Resources for Types of Triangles Once I complete notes with my students, I want them to be applying what they have learned. If you don’t use it, you lose it, right? When I am planning practice activities, I want to make sure they are interactive and allow everyone to have a seat at the table. Math anxiety is real! It’s important to be aware of all the emotions your kiddos may be feeling when it comes to math. With that in mind, I have included some types of triangle activities that have students practicing the skills in a collaborative, low-intensity manner. Classifying Triangles Footloose Math Task Cards Activity and Math Wheel An awesome classroom activity I use with my students to review types of triangles and to get them up and moving is my Types of Triangles Footloose cards. Footloose isn’t just your typical sit-and-listen routine. I hand out these cards, and suddenly, we’re not just classifying triangles by sides and angles. We take it a step further by finding those missing angles and sides like geometry detectives. Sometimes, alarms go off in my students’ heads when they discover a missing angle. We pull out our math doodle wheels and remind ourselves that all three angles within a triangle equal the grand number of 180! I tend to shake up the delivery of the task cards because sometimes I have them posted around the classroom ahead of time. Then they travel around to each card with a partner to solve and record answers on their paper. With that being said, Footloose isn’t a one-size-fits-all deal. You can roll it out for the whole class, turn your class into an interactive experience, use it in small groups, or throw it into center rotations. It’s flexible, it’s fun, and it’s got the whole class moving! Triangles Color by Number Activity Missing Angles and Sides My Triangles Color by Number activity is a low-intensity math activity for different types of angles and missing angles with a dash of color. Students solve problems while turning their answers into My students tackle 20 problems that dive deep into triangle types. Problems range from finding missing side lengths and angle measures to calculating the perimeter of triangles. Instead of the usual pencil and paper routine, they find their solutions on a coloring sheet/section. I love using this resource for many reasons, but one of my favorites is that it’s self-checking! If their answer isn’t on the coloring sheet, they know it’s time for a quick math check. It’s instant feedback that we know and love! Plus, there are both print and digital versions available. I like flexibility, which is the name of the game for us teachers, so I wanted this resource to be flexible for your teaching style. Classifying Triangles Math Game | Math Activity | Truth or Dare Truth or Dare? I choose truth! I’ve got a game-changer for your next class or review session for identifying types of triangles that’s as fun as it is educational. I’ve crafted 36 question cards that’ll turn your class into a buzz of excitement. It’s not your average Q&A session because this one comes with a “truth” or “dare” twist! I pull this Truth or Dare activity out for center time when I am working with small groups, will be out of the classroom, or for review sessions. My students love it when they see this activity on the schedule. Your students dive into a world where they get to choose their fate. Questions focus on the different triangle types (acute, obtuse, right, isosceles, equilateral, scalene) and essential concepts like the 180 degrees in a triangle. They earn a 1-point “truth” for those classic true or false questions. If they’re feeling daring, there are “dare” cards worth 2 or 3 points, throwing in some challenges. Think along the lines of finding missing angles or sides, explaining why a triangle is classified a certain way, or even determining if a shape can make the triangle cut based on angle measures or side lengths. The beauty of this is it’s not just a solo gig. Gather your kiddos in groups of 3-4, and let the game begin. Keep in mind the smaller groups mean more questions tackled by each student, meaning more Time to Teach Types of Triangles Exploring types of triangles with the math doodle wheel makes note-taking into a masterpiece, it’s all about transforming math anxiety into math excitement. These resources are tools that keep on giving throughout the unit and year. The adventure continues with activities that have your mathematicians applying their skills through interactive activities like Footloose, Color by Number, and Truth or Dare. Get ready to witness “Aha!” moments, foster collaborative learning, and make math a destination worth exploring. Save for Later Remember to save this post to your favorite math Pinterest board so you can quickly come back when you are planning your types of triangles lessons!
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Unknown in Multiplication (examples, solutions, videos, worksheets, homework, lesson plans) Related Topics: Lesson Plans and Worksheets for Grade 3 Lesson Plans and Worksheets for all Grades Common Core For Grade 3 More Lessons for Grade 3 Math Examples, videos, and solutions to help Grade 3 students learn how to interpret the unknown in multiplication and division to model and solve problems. Common Core Standards: 3.OA.3, 3.OA.4, 3.OA.5, 3.OA.7, 3.OA.1, 3.OA.2, 3.OA.6, 3.OA.8 New York State Common Core Math Grade 3, Module 3, Lesson 7 Worksheets for Grade 3 Concept Development Thad sees 7 beetles when he weeds his garden. Each beetle has 6 legs. How many legs are there on all 7 beetles? Lesson 7 Homework 1. Match the words on the arrow to the correct equation on the target. 2. Ari sells 6 boxes of pens at the school store. a. Each box of pens sells for $7. Draw a tape diagram and label the total amount of money he makes as m. Write an equation and solve for m. b. Each box contains 6 pens. Draw a tape diagram and label the total number of pens as p. Write an equation and solve for p. 3. Mr. Lucas divides 28 students into 7 equal groups for a project. Draw a tape diagram and label the number of students in each group as n. Write an equation and solve for n. Try the free Mathway calculator and problem solver below to practice various math topics. Try the given examples, or type in your own problem and check your answer with the step-by-step explanations. We welcome your feedback, comments and questions about this site or page. Please submit your feedback or enquiries via our Feedback page.
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Cayley-Bacharach applications The Cayley-Bacharach theorem in projective geometry states that if two cubics share nine points (in sufficiently general position) and a third cubic passes through eight of those points, then it also passes through the ninth. Here, ‘sufficiently general position’ means that as long as you’re not deliberately stupid with applying this theorem, it will work. I’ve heard this referred to as a ‘super-theorem’, because its main use is to prove simpler theorems. Firstly, note this very important fact: • The union of a degree-m and a degree-n algebraic curve is a degree-(n+m) algebraic curve. So, for cubics, we can have either an irreducible cubic such as an elliptic curve, or the union of a conic and a line, or indeed the union of three lines. In most applications, we don’t use the full force of Cayley-Bacharach, and concentrate instead on these ‘degenerate’ cases. Pascal’s theorem (and converse) Pascal’s theorem states that if we inscribe a hexagon in a conic, the extensions of opposite sides intersect at three collinear points. The hexagon can self-intersect, such as the following example: I’ve coloured each cubic in a different colour, such that the nine points are triple intersections. If we lack a single piece of information (such as a concurrency or collinearity) but know that all the other pieces of information hold, we can deduce the final piece of information through Cayley-Bacharach. For instance, the statement of Pascal’s theorem says that the red line exists, and the converse states that the red conic exists. To determine on which branch of the cubic the point lies (either the conic or the line), we can just apply Bezout’s theorem. Projective dual: Brianchon’s theorem Once we have Pascal’s theorem, we can take poles and polars about an arbitrary conic to obtain the projective dual theorem. In projective duality, we interchange lines and points, whereas conics remain conics. Concurrency of lines corresponds to collinearity of points, and so on. The projective dual of Pascal’s theorem is called Brianchon’s theorem, which states that a hexagon circumscribed about a conic has concurrent diagonals. I’ve colour-coordinated this diagram to correspond with the diagram of Pascal’s theorem; the black lines were originally black points, and the coloured points were originally lines of that colour. We can tell that this is a projective dual of a C-B corollary rather than a C-B corollary itself, because the lines are black and the points are colourful. Both Pascal’s theorem and Brianchon’s theorem have special ‘degenerate’ cases. If the conic in Pascal’s theorem degenerates into a pair of lines, we obtain Pappus’s theorem instead. This is also a direct corollary of Cayley-Bacharach, for obvious reasons. In the same way that Pappus is a degenerate case of Pascal, we can view Pascal as a degenerate case of the theorem obtained by replacing the red conic and line with a single elliptic curve. This associativity law is used in elliptic curve cryptography, and more recreationally in my elliptic curve calculator. It is the most difficult part of showing that the points on an elliptic curve form a group under a geometric operation. Radical-axis theorem and pivot theorem When we apply Cayley-Bacharach, we actually use the complex projective plane rather than the ordinary Euclidean plane. There are two invisible ‘circular points at infinity’, through which all circles pass. Using these, we can derive some theorems involving circles from Cayley-Bacharach. Only seven of the points are visible here; the other two are the circular points. With three circles and three lines, there are two distinct theorems. Firstly, the pivot theorem: We also obtain the radical-axis theorem in this way: Circle inversion When we’re dealing with circles and lines rather than arbitrary conics, we can concentrate on the real Euclidean plane and extend it to give the Riemann sphere, or complex projective line. This is quite a sneaky trick to switch seamlessly between these two universes, but we can get away with it here. Inverting the pivot theorem yields Miquel’s theorem. This states that if we have eight points and ‘faces’ correspond to concyclic points, then five faces of a cube determine the sixth face: The diagram above was obtained by stereographic projection, so that the cubic nature can easily be appreciated. The radical axis theorem treats the six ‘diagonal’ planes inside a cube instead of the six faces. In the very rare instances where both apply simultaneously, we get the famous orthocentric configuration. Quartic Cayley-Bacharach It turns out that there is a generalisation of Miquel’s theorem to eight conics and sixteen points, but this requires the quartic analogue of Cayley-Bacharach to prove. I’ll leave this particular one as an exercise to the reader, and instead use QCB to derive some other beautiful results. There was a particularly elegant problem posted on Art of Problem Solving, the solution to which I’ll describe here. Firstly, it’s necessary to know the statement of quartic Cayley-Bacharch. If we have 16 sufficiently-general-position points lying on two quartics, and a third quartic passes through 13 of those points, then it also passes through the other 3. In particular, if we need to prove that eight points to lie on a conic, five of them automatically do and we can get the other three through Cayley-Bacharach. Let’s have a look at the AoPS problem: We want to show that the centres of the eight pink/purple circles lie on a conic. We’ll apply the fact that they lie on the angle bisectors of the triangles. Indeed, once we’ve drawn the angle bisectors, we can forget about the grey lines and pink/purple circles completely: The red and blue angle bisectors intersect to give the eight incentres. However, they also appear to converge on the circumference of the circle to give another four points. This is not a coincidence, and can be deduced from the converse of ‘equal arcs subtend equal angles’. These sixteen points lie on the blue quartic and the red quartic, and we can also draw a green conic through five of the remaining eight points. Quartic Cayley-Bacharach then reveals that the last three points also lie on this conic: This diagram is a generalisation of Pascal’s theorem, which is a corollary of generalised Cayley-Bacharach. In particular, if we have a 2n-gon inscribed in a conic, and colour alternate edges red and blue, the remaining n^2 − 2n points lie on an algebraic curve of degree n − 2. This is explored in a paper by Gabriel Katz. Solving an AMS problem We’ll solve a non-trivial olympiad-style problem using as much pure projective geometry as possible. This differs completely from any ordinary way of solving the problem, although it’s somewhat neater (the standard solution has horrible Menelaus bashes and so forth). Notice that Q and R are defined in quite a contrived way, so we’ll see if we can learn more about these points. It’s always good to draw a diagram in these contexts; I’ll leave it relatively uncluttered and just include F, E, B, C, Q and R (these are the only ones we’ll require at this stage). FE and QR meet at a point on the line at infinity, which is thus collinear with the two circular points. Colour the line at infinity red. We already have that the green circle exists (and has diameter BC, by Thales’ theorem), so we can conclude that BCRQ lie on a blue circle by Cayley-Bacharach. You could alternatively use an angle chase, but that’s too boring for cp4space. The orthocentre exists; hence, (P,D;B,C) is a harmonic range. Consequently, the circle on diameter PM is the inverse of the line AD (polar of P) with respect to the circle on diameter BC (centre M). So, AD is the radical axis of circles on diameters PM and BC. This means we can just apply a degenerate case of the radical axis theorem to the points {D, Q, R, B, C, P, M} to deduce that {M, P, Q, R} are indeed concyclic. 7 Responses to Cayley-Bacharach applications 1. That last paragraph was hard to follow – it would have been nice if you’d included another diagram. 2. Heello mate great blog post This entry was posted in Uncategorized. Bookmark the permalink.
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Does Sex-Differential Gene Expression Drive Sex-Differential Selection in Humans? Sex differences in human transcriptomes have been argued to drive sex-differential natural selection (SDS). Here, we show that previous evidence supporting this hypothesis has been largely unfounded. We develop a new method to test for a genome-wide relationship between sex differences in expression and selection on expression-influencing alleles (eQTLs). We apply it across 39 human tissues and find no evidence for a general relationship. We offer possible explanations for the lack of evidence, including that it is due in part to eQTL ascertainment bias towards sites under weak selection. We conclude that the drivers of ongoing SDS in humans remain to be identified. Sex-differences in gene expression have been theorized to be the result of long-term sex-differential selection (SDS), where allelic fitness effects differ between males and females^1,2. When a gene product is only beneficial in one sex, it is expected that expression modifiers will evolve to increase expression for that sex and decrease expression in the other ^3. The reverse causality has also been proposed: sex differences in gene expression might be the driver of SDS on expression modifiers^1,4. SDS acting on viability within the current generation can generate between-sex differences in allele frequency^5,6 at expression modifiers (expression Quantitative Trait Loci, or eQTLs)^3. Sexually dimorphic traits (diseases being one example with strong repercussions on survival^6,7,8,9) may be subject to SDS acting through viability. Previous work has proposed and tested theoretical models for relating divergence in eQTL allele frequencies with differential expression between sexes^1,10,11. Some of the empirical results have, however, been called into question^12,13,14. A controversial result by Cheng and Kirkpatrick (2016)^1 (hereafter “CK16”) is a characteristic pattern relating between-sex F[ST]^1,5,15 (henceforth “F[ST]”) to sex differences in gene expression (henceforth “Δ”). They observed high F[ST] values at intermediate values of Δ and near-zero values of F[ST] when a gene is expressed evenly between the sexes (Δ = 0) or is only expressed in one sex (Δ = ±1). They nicknamed this bimodal pattern “Twin Peaks”. The Twin Peaks pattern has been used as a signal to detect SDS in several species^1,4,16,17. Here, we revisit the CK16 model and statistical approach. We find that their interpretation was based on previously unappreciated statistical artifacts. We then refine their model and apply it to new, more extensive data on gene expression and allele frequencies. Across 39 human tissues, we find no evidence for a genome-wide relationship between viability SDS and Δ. We discuss how a bias in eQTL discovery towards eQTLs under weaker selection can explain the lack of signal for a relationship between SDS and Δ, and how the drivers of SDS may still be investigated. Twin Peaks is a Statistical Artifact We begin by revisiting the Twin Peaks pattern with a critical eye towards caveats in its application and interpretation. Most importantly, we reconsider the model that CK16 proposed to explain the pattern. It is based on two key assumptions. First, the relationship between a gene’s expression levels and its effect on fitness in each sex is linear. Second, sexually antagonistic selection is symmetric, meaning selection coefficients are equal in magnitude and opposite in sign between sexes. This yields a quadratic relationship between F[ST] and Δ at a biallelic site affecting expression. At small absolute values of Δ, where p is the allele frequency in zygotes and q = 1 − p. F[ST] is the between-sex fixation index^18 which is used to quantify allele frequency differences between males and females. In the model, these differences are due solely to the sex differences in post-zygotic fitness effects of alleles. The quantity Δ is the sex difference in gene expression (Methods). Finally, A is a compound parameter involving the within-sex effect of gene expression level on fitness. The expectation in CK16 for F[ST] at extreme expression differences (|Δ| → 1), however, is based on intuition rather than a model. The authors suggest that if a gene is not expressed in one sex (Δ = ±1), then selection will not act on it. Selection on the other sex should optimize expression levels, so under the symmetrical selection assumption there will be no ongoing directional selection in that sex either. As neither sex experiences selection, there will be no force driving increased F[ST]. CK16 then interpolate a bimodal shape by joining the quadratic relationship at low |Δ| and the expectation for F[ST] = 0 at Δ = ±1 (Fig. 1). However, the assumption of symmetric selection used at low values of Δ may be inappropriate for large Δ values. In particular, when Δ = ±1, the lack of expression in one sex plausibly suggests different selection between males and females. It is therefore not intuitive that F[ST] should simply go to zero at sites regulating expression in these genes. Additionally, although the quantity 2pq appears in Eq. 1, CK16 did not include that term in fitting the model, effectively assuming it to only contribute random noise to the relationship between F[ST] and Δ. Because of these caveats to the model, we decided to revisit the support for Twin Peaks. We first ask whether the Twin Peaks pattern is due to SDS by applying the statistical tests of CK16 to data generated under a null hypothesis of no relationship between Δ and SDS. We replicated the pattern shown in CK16 by using the same data and methods. Namely, we performed a 4^th degree polynomial regression of F[ST] on Δ, using allele count data from 1000 Genomes^19 and expression data from the gonads (ovaries and testes) in GTEx v3^20. The curve and associations generated using these datasets we refer to as the “real data” (blue line in Fig. 2a). We then generated an empirical null by permuting ovary and testes tissue labels in the expression data but retaining sex labels associated with F[ST] values (Methods), then recomputed Δ values, and again fit a 4^th-degree polynomial. We find 21% of permutations qualify as Twin Peaks according to the criteria used by CK16, and that the polynomial fit to the real data is not visually distinct from those fits to the null (Fig. 2a). Higher order polynomial regressions can also yield spurious fits because distant points have an oversized impact^21,22. We therefore binned genes by Δ values and examined the relationship with mean F[ST] in each bin. Again, the real data shows no unusual relationship between F[ST] and Δ compared to null data (Fig. 2b). From these results, we conclude the Twin Peaks pattern is not statistically significant. Importantly, this permutation method differs from the method used in CK16, which permuted Δ values across genes. Both methods break the associations between F[ST] and Δ as desired. Our sex -label permutation method, however, preserves F[ST] associations with the gene’s overall expression, maintaining gene features such as expression variance which the CK16 method does not. This explains why their reported p-value for Twin Peaks curves (0.016) is lower than our (0.21). We hypothesized that one reason for the spurious Twin Peaks pattern is confounding. In particular, variance in expression (regardless of sex) is positively correlated with both Δ and F[ST]. Previous work has shown that genes subject to weaker stabilizing selection show higher variance in expression^23. Higher variance means larger differences between randomly selected subgroups, and therefore it should translate to larger values of Δ even in the absence of SDS. In turn, stronger (sex-agnostic) selection can lead to stronger drift at linked sites^24,25,26,27. Taken together, the confounding with variance in gene expression could generate a relationship between SDS and |Δ| in the absence of SDS. Indeed, expression variance and F[ST] are positively correlated (Pearson p = 0.011; Fig. 2c), consistent with this hypothesis. In sum, we do not find support for Cheng and Kirkpatrick’s conclusion that there is a genome-wide relationship between F[ST] and Δ. No evidence for genome-wide SDS on eQTLs Although we do not find support for a relationship between sex differences in gene expression and selection using CK16’s methodology for generating Twin Peaks, one may still exist. We test this hypothesis across many tissues using improved statistical modeling, data, and methods. Despite the caveats to portions of the model discussed above, we believe that CK16’s theoretical expectation of a quadratic relationship between F[ST] and small values of Δ is valid. We therefore built on that model by introducing the compound parameter δ^2 = 4qpΔ^2 and rewriting Eq. 1 as (Methods). To estimate A, for each gene-tissue pair, we used the single cis-eQTL from GTEx v8^28 with the strongest association with its expression. Our estimator of A is then the inverse-variance-weighted linear regression of F[ST] on δ^2 (Eq. 2). The advantages of this formulation over CK16’s are that it allows direct estimation of A, the strength of SDS on sex differences in expression. This model also accounts for variation in allele frequencies across sites. Further, by using a single eQTL to calculate F[ST] for the whole gene, we circumvent biases which can arise when using a simple mean to estimate gene-wide F[ST]^29,30. To accompany the updated regression methodology, we also updated the datasets for both allele frequencies and gene expression. For allele frequencies, we used the Non-Finnish European subset in gnomAD v3 (averaging over 60,000 allele samples per site)^31. We used expression data from 49 tissues from the GTEx v8 dataset^28 (averaging over 200 male samples and 150 female samples). Both datasets greatly expand our sample size compared to CK16, and GTEx v8 provides sample-specific sex labels for calculating Δ across multiple tissues instead of just the gonads—ovaries and testes—as in the original Twin Peaks paper. Using the updated statistical framework method, we find no evidence that A differs significantly from zero in any tissue (Fig. 3; Methods). One explanation for the absence of a pattern may be found in recent work by Mostafavi et al. (2023). The authors contrasted how selection impacts the discovery of genome-wide association study (GWAS) hits with the discovery of eQTLs. Variants with large effects on phenotypes are expected to segregate at low frequencies, reducing discovery power. In GWAS, this is counterbalanced by increased power due to their large effect sizes. Consequently, in GWAS, low frequency variants can still be detected if their effect is large enough. In contrast, eQTL discovery is based only on the effect of genotype on gene expression which does not necessarily translate to fitness-relevant trait variation. Detection of strongly selected sites is therefore less likely^32. This ascertainment bias can weaken the relationship between Δ and F[ST] at eQTLs, as sites with high F[ST] are less likely to be eQTLs. In “Twin Peaks is a Statistical Artifact”, we suggested that large sex differences in expression may be entirely unrelated to SDS—for instance, merely tagging genes with high expression variance. There are other potential explanations for the lack of a genome-wide relationship. While some sex differences in expression may be driven by or drive SDS, these are the exception rather than the rule. The majority of sex differences in expression may be a regulatory side effect of SDS on different genes. Lastly, sex-differential gene expression may be due to past SDS, but not drivers of current SDS^33. Regardless of the reason, if a causal relationship between Δ and FST is rare in the genome, it would be difficult to detect using models that assume a pervasive, persistent relationship between the two. Previous work suggested a genome-wide relationship between sex differences in expression and SDS. However, this work was based on statistical artifacts and confounded effects, such as stabilizing selection and variance in gene expression. Even when using newer data and improved statistical methods, we found no evidence for a genome-wide relationship between sex-differential expression and contemporary SDS. In contrast, studies that measured SDS irrespective of expression or trait variation have reported pervasive, genome-wide signals of SDS in the human genome^5,10,34. While causal relationships, past and present, between SDS and sex-differential gene expression in humans remain plausible, they are yet to be fully elucidated. Materials and Methods Permuting Twin Peaks sex-labels In the section “Twin Peaks is a Statistical Artifact”, we demonstrate that the Twin Peaks curve presented in CK16 is not statistically significant compared to a permuted null. To do so, we used the methods and datasets described in CK16 and Eq. 1. We computed allele count using the 1000 Genomes dataset^19 and calculating F[ST]. We filtered out any sites with only a single alternative allele in either males or females (i.e., singletons). We used the Transcripts Per Million (TPM) normalized GTEx v3 (referred to as “pilot” on the download page) dataset^20 for expression levels for calculating Δ. Because GTEx v3 does not have individual sample labels, this analysis only compared expression in the gonads (ovaries and testes) where expression level summaries represent a single biological sex. We used the Ensembl GRCh38 v77^35 annotation file for gene annotations. Following CK16, we limited our analysis to protein-coding genes. To estimate F[ST], we used Hudson’s estimator^36 based on the R package used in CK16^37 as presented by Bhatia et al. (2013)^29, Here, pm and pf are the allele frequencies in males and females respectively, and nm and nf are the number of males and females respectively. As a measure of sex-differential expression we used Δ as defined in CK16: Here, x[m] and x[f] are sex-averaged TPM-normalized expression levels in males and females respectively. Because in Eq. 3 is a site-specific estimator while Δ is gene-wide, we generated a gene-wide estimate of F[ST] by taking a simple average of all sites within a gene body plus 1000bp upstream and downstream. To generate the Twin Peaks curve, we used a 4^th-degree polynomial regression between and Δ. We note that here and in CK16, this regression therefore ignores variation in heterozygosity across sites (Eq. 1). In our improved model (Results: “No evidence for SDS on eQTLs across all genes in multiple tissues” and Methods: “Applying a new model and method for SDS-expression regression”) we correct this omission. To compare the original Twin Peaks curve to curves generated under the null of no SDS, we generated an empirical null distribution of 4^th-degree regression curves using sex-label permutations of gene expression data. We permuted the ovary and testes labels for each GTEx sample, then recalculated Δ for all genes. We then re-performed the 4^th-degree polynomial regression on the new Δ values (the F[ST] values remain unchanged). This was repeated 500 times. By permuting sex labels, we break the association between sex-differential expression and sex-differential expression, while preserving other gene-level features. To quantitatively evaluate the significance of Twin Peaks in the null distribution, we used the three criteria laid out by Cheng and Kirkpatrick (2016) for classifying a curve as Twin Peaks. Namely, a 4^th-degree polynomial must 1) be significant (p < 0.05) for the 4^th-degree term by ANOVA, 2) have a negative coefficient for the Δ^4 term, and 3) have three real roots. Any 4^th-degree regression passing all three criteria is classified as Twin Peaks. Applying a new model and method for SDS-expression regression In the section “No evidence for SDS on eQTLs across all genes in multiple tissues”, we described a revised method for testing a genome-wide relationship between SDS and Δ. For this analysis, we revised the inference model and data. We used the Non-Finnish European subset in the gnomAD V3 dataset^31 to calculate between-sex F[ST] and heterozygosity at each eQTL. This set of samples has an order of magnitude more individuals (average of 31,470 samples per site) compared to 1000 Genomes data (average of 2,300 samples per site when combining all ancestry groups) originally used in CK16. Additionally, we used the normalized TPM from the current GTEx v8 for expression level data, which includes many more samples (average of 236 male and 168 female samples per tissue) than the pilot (8 male samples and 13 female samples in the gonads)^28. This version also provides sample-specific sex labels, enabling us to use additional tissues beyond gonads. The equation we base our regression on is the model relating F[ST] with δ^2 shown in Eq. 2. Now, rather than performing a 4^th-degree polynomial regression as in the analysis for “Twin Peaks is a Statistical Artifact”, we perform a weighted linear regression of F[ST] to δ^2 = 4pqΔ^2, isolating A as the coefficient of regression. Note that by including 4pq in the independent variable, we allow information about heterozygosity to affect the SDS-expression relationship. We weighted each point of the regression by 1/Var(expression), where Var(expression) is the averaged variance in expression of each sex. This should decrease the leverage of points with large |Δ|. To get a gene-wide estimate of F[ST], we used eQTLs from GTEx as mapped in the v8 study^31. For each gene, we chose the single cis-eQTL (within 1Mbp of the gene’s midpoint) with lowest p-value for association with the gene and use that site for calculating F[ST], p, and q. By using an eQTL selected this way, we tried to isolate the effect of selection on gene expression to the site presumed to be contributing most to expression changes in that gene. To determine the significance of our A estimates, we generated 90% null acceptance regions by permuting sex labels. We permuted sex labels as in “Twin Peaks sex-label permutation” such that for each tissue, sex labels are permuted in GTEx expression samples and Δ is recalculated across all genes. Then, for each iteration A is recalculated by linear regression using the new Δ values (F[ST] and 4 pq remain unchanged). The 90% null acceptance region was obtained by the 5^th and 95^th percentiles of 1,000 A values calculated on permuted expression data. We thank Jared Cole for comments on the manuscript and other members of the Harpak Lab for helpful conversations. The work was funded by NIH grants GM11685307 to M.K., and R35GM151108 and a Pew Scholarship to A.H.
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Energy Equation Next: Polytropic Relation Up: Basic Equation of Fluid Previous: Expression for Momentum Density &nbsp Contents The above basic equations (A.3) and (A.10) or equations (A.7) and (A.11) contain three dependent variables v. The number of the variables, 3, is larger than the number of equations, 2. Therefore, an extra equation is needed to close the basic equations. Kohji Tomisaka 2007-07-08
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While going through my puzzle files, I came across the Slovak Championships folder. I realised I had never shared these puzzles on my blog. So I figured I might as well share them now. I was contacted by Zuzanna Hromcova to write puzzles for their championship. We were given a number of categories to write puzzles in. One of the categories was non-grid puzzles, for which I provided three different genres; namely ABC Decoder, Dice and Mastermind. Dice and ABC Decoder are types I enjoy writing; Mastermind I didn't have that much experience with. But it was something I'd like to give a go. The other categories I picked were Latin Squares and Division puzzles. For each type we had to write a standard genre and a variant on the genre. I picked Skyscrapers, with Haido as the variant. I like Skyscrapers and I thought Haido still had part of the Skyscraper logic, but used differently enough to make it not like solving four skyscraper puzzles. For the Division set I picked ABCD Division, with Sum Division as the variant. It's a type I have seen a lot when I first started puzzling, but I haven't really seen it much since. I thought sums was an obvious variant, but I haven't really seen it this way much. I have seen a similar variant where the grid has to be divided into a complete set of pentominoes, but not really without this I tried to put a bit of theming in the non-grid puzzles. I wrote a few nine digit ABC Decoders for the 2014 24 hour championships, and I thought that was a good size to use in a championship. The letters spell out THE SLOVAK, which was the nicest way I could use nine different letters to write something Slovakia related. I found some words with opposite meanings in the letter set, so I used those. I think it turned out well. I used a similar opposites theme for the Dice puzzle, with an addition of 5 words to make it unique. I think not all words are necessary for uniqueness, but it solves pretty well this way. The first Mastermind puzzle looks really nice, with a sequence of numbers and only white circles it solves really nicely. The second one was merely an attempt to construct a nice logical 5 digit I thought both Skyscrapers puzzles turned out nicely. The first puzzle uses three 4s and three 5s. The second puzzle has a trio of the same digit on each side. Of course I couldn't use four different digits as these are the only three digits you can have three of the same clue on the same side in this size. I find it hard to theme Haido puzzles as the clues are a bit limited, but they both have nice logical paths. The first time I saw an ABCD puzzle this way was at a Dutch championship. It was a bit of a surprise then. I wrote a similar puzzle for puzzlepicnic once and I thought it would be fun to include one for the championship. The ABCDE puzzle is a standard layout and I think it solves well. The sum puzzles were a bit hard to work out openings at first as there are so many ways to reach the sums. So I went with obvious opening digits for both puzzles to then work back to more ambiguous digits towards the end. I think they both turned out well. Puzzles can be found below. Read more »Posted in Dice, Division, Haido, Mastermind, Puzzle, Puzzle Championship, Skyscrapers, Word
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Cubic Inches to Liters Conversion (in³ to L) - Inch Calculator Cubic Inches to Liters Converter Enter the volume in cubic inches below to convert it to liters. Do you want to convert liters to cubic inches? How to Convert Cubic Inches to Liters To convert a measurement in cubic inches to a measurement in liters, multiply the volume by the following conversion ratio: 0.016387 liters/cubic inch. Since one cubic inch is equal to 0.016387 liters, you can use this simple formula to convert: liters = cubic inches × 0.016387 The volume in liters is equal to the volume in cubic inches multiplied by 0.016387. For example, here's how to convert 5 cubic inches to liters using the formula above. liters = (5 in³ × 0.016387) = 0.081935 L Cubic inches and liters are both units used to measure volume. Keep reading to learn more about each unit of measure. What Is a Cubic Inch? A cubic inch is a unit of volume equal to the space consumed by a cube with sides that are one inch in all directions. One cubic inch is equivalent to about 16.387 cubic centimeters or 0.554 fluid The cubic inch is a US customary and imperial unit of volume. A cubic inch is sometimes also referred to as a cubic in. Cubic inches can be abbreviated as in³, and are also sometimes abbreviated as cu inch, cu in, or CI. For example, 1 cubic inch can be written as 1 in³, 1 cu inch, 1 cu in, or 1 CI. Learn more about cubic inches. What Is a Liter? A liter is a unit of volume equal to 1,000 cubic centimeters or 0.264172 US gallons.^[2] The liter is a special name defined for the cubic decimeter and is exactly equal to the volume of one cubic decimeter (1 decimeter is 1/10 of a meter, or 10 centimeters). The liter is an SI accepted unit for volume for use with the metric system. A liter is sometimes also referred to as a litre. Liters can be abbreviated as L, and are also sometimes abbreviated as l or ℓ. For example, 1 liter can be written as 1 L, 1 l, or 1 ℓ. Learn more about liters. Cubic Inch to Liter Conversion Table Table showing various cubic inch measurements converted to liters. Cubic Inches Liters 1 in³ 0.016387 L 2 in³ 0.032774 L 3 in³ 0.049161 L 4 in³ 0.065548 L 5 in³ 0.081935 L 6 in³ 0.098322 L 7 in³ 0.114709 L 8 in³ 0.131097 L 9 in³ 0.147484 L 10 in³ 0.163871 L 11 in³ 0.180258 L 12 in³ 0.196645 L 13 in³ 0.213032 L 14 in³ 0.229419 L 15 in³ 0.245806 L 16 in³ 0.262193 L 17 in³ 0.27858 L 18 in³ 0.294967 L 19 in³ 0.311354 L 20 in³ 0.327741 L 21 in³ 0.344128 L 22 in³ 0.360515 L 23 in³ 0.376902 L 24 in³ 0.39329 L 25 in³ 0.409677 L 26 in³ 0.426064 L 27 in³ 0.442451 L 28 in³ 0.458838 L 29 in³ 0.475225 L 30 in³ 0.491612 L 31 in³ 0.507999 L 32 in³ 0.524386 L 33 in³ 0.540773 L 34 in³ 0.55716 L 35 in³ 0.573547 L 36 in³ 0.589934 L 37 in³ 0.606321 L 38 in³ 0.622708 L 39 in³ 0.639095 L 40 in³ 0.655483 L 1. National Institute of Standards and Technology, Specifications, Tolerances, and Other Technical Requirements for Weighing and Measuring Devices, Handbook 44 - 2019 Edition, https:// 2. National Institute of Standards and Technology, Units outside the SI, https://physics.nist.gov/cuu/Units/outside.html More Cubic Inch & Liter Conversions
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"IT Organization" name in WinPE Can someone please point me in the right direction so I can set the "IT Organization" name in the Task Sequence progress window. I've set this in the custom client settings which is deployed to the same collection as the task sequence. When I run a task sequence that captures User State (requires task sequence to be launched from Windows) the "IT Organization" works fine and it displays our company name. However, when I PXE boot a task sequence for a task sequence with no USMT capture, it only displays "IT Organization" in the task sequence progress window. This is when I launch from Software Center: it gets the name from the Computer Agent part of the Default Client Settings defined for the hierarchy, • 1 I face the same, I am 100 % sure that I have custom client settings and deployed this to target collection (not been deployed to UNKNOWN). But still it says running IT Organization. Created a client settings again, deployed it to a different collection, same results. Hi Peter, Thanks for reply. I can update the DP (which I have tested already) but cant Re-Distribute since there is no DP available if I try to distribute the contents again. Is it something I am missing? Hi Niall Yes, I am updating DP. I followed your famous sccm guides in LAB. Created a couple of custom client device settings, deployed them to specific test collections, created TS and checked, all gives running IT Organization. Today, I changed computer agent settings in DEFAULT client settings, now I can see that it appears the settings I provided in DEFAULT client settings. Note: I have multiple boot images, one set is default x64 , x86 and other is copy of default boot images but I have changed background /PE wallpaper. Any thoughts? do as follows, change your organization name in Default Client settings, once done, update the boot image that is linked in your task sequence to it's distribution points, that should solve it. to identify the boot image, right click on the task sequence and you'll see it listed under advanced.
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5 Corrected constraint logic programming exercises - Complex systems and AI 5 Corrected Constraint Logic Programming Exercises Consider the following puzzle: Each region of the grid must be filled with a number between 0 and 9 so that • the numbers in two adjacent regions (vertically or horizontally) are different • whenever there are four regions that meet at a point (indicated by a small circle), the sum of their numbers is equal to 20. Solve this problem by constraint logic programming. Let's name from top left to bottom right (by line) the empty squares by a letter. Each domain has a value between 0 and 9. Thereafter, the two constraints must be filled in manually. For each square, its score must be different from its neighbours. For each circle, the sum must be equal to 20. We therefore have the following constraint logic programming: L = [A,B,C,D,E,F,G,H,I,J,K,M], At #\= 7, At #\= 8, At #\= 6, B #\= 8, B #\= 4, C #\= 7, C #\= 6, C #\= 9, D #\= 6, D #\= 8, E #\= 6, E #\= 8, E #\= 4, E #\= 5, F #\= 9, F #\= 6, F #\= G, F #\= Y, G #\= 6, G #\= 5, G #\= K, H #\= 5, H #\= 4, H #\= M, I #\= 9, I #\= J, J #\= K, K #\= 5, K #\= M, 7+6+A+C #= 20, A+6+8+D #= 20, 6+8+D+E #= 20, 4+8+B+E #= 20, 9+6+C+F #= 20, 5+6+E+G #= 20, 4+5+E+H #= 20, 9+F+I+J #= 20, F+G+J+K #= 20, 5+H+K+M #= 20, Consider the game Bokkusu. The object of the game is to mark certain squares of the given grid in black. Each box has two values: an A value and a B value. The numbers on the right denote the A values of the boxes in each corresponding row and the numbers at the bottom denote the B values of the boxes in each corresponding column (These values always range from 1 to the size of the grid by increasing by 1 from left to right or from top to bottom). The values on the left indicate the sum of the B values of the black boxes for each row and the values on the top indicate the sum of the A values of the black boxes for each column. An example of a grid and its solution is given in the figure. In the solution of this example we have for example 7 = 1 + 2 + 4 (the boxes with B value of 1, 2 and 4 are marked black) and 6 = 2 + 4 (the boxes with A values of 2 and 4 are marked). Solve this problem by constraint logic programming. Let's name each square in the grid in a matrix fashion. For the constraints, the sum of the binary variables (white=0 or black=1) multiplied by the coefficients on the left or on the top must be equal to the value opposite. L = [X11,X12,X13,X14, X11 + 2*X12 + 3*X13 + 4*X14 #= 3, X21 + 2*X22 + 3*X23 + 4*X24 #= 7, X31 + 2*X32 + 3*X33 + 4*X34 #= 1, X41 + 2*X42 + 3*X43 + 4*X44 #= 2, X11 + 2*X21 + 3*X31 + 4*X41 #= 5, X12 + 2*X22 + 3*X32 + 4*X42 #= 6, X13 + 2*X23 + 3*X33 + 4*X43 #= 1, X14 + 2*X24 + 3*X34 + 4*X44 #=2 We consider the problem of generating a timetable. There are 7 classes and 4 slots. A slot can accommodate a maximum of 2 lessons. There are the following constraints: – Course 4 must be before course 6. – Course 5 must be before course 7. – Course 6 must be before course 2. – Course 1 cannot be in parallel with courses 2, 3, 4 and 7. – Course 2 cannot be in parallel with courses 3 and 6. – Course 3 cannot be in parallel with courses 4, 5 and 6. – Course 4 cannot be in parallel with courses 5 and 6. – Course 5 cannot be in parallel with course 7. Solve this problem by constraint logic programming. We will consider that the courses are domains, and that they can take a value from 1 to 4 corresponding to the niche that will be chosen. Constraint logic programming is as follows: L = [C1,C2,C3,C4,C5,C6,C7], C4 #< C6, C5 #< C7, C6 #< C2, C1 #\= C2, C1 #\= C3, C1 #\= C4, C1 #\= C7, C2 #\= C3, C2 #\= C6, C3 #\= C4, C3 #\= C5, C3 #\= C6, C4 #\= C5, C4 #\= C6, C5 #\= C7, Inshi no heya is a Japanese logic game. It is played on a rectangular grid of cells which are separated into rectangles. One dimension of each rectangle is 1 while the other dimension varies. Each rectangle is placed horizontally or vertically and contains a number. The goal is to fill all cells with numbers from 1 to 9 so that: – If we multiply all the digits of each rectangle we obtain the number indicated in the rectangle – No number appears more than once in a column. – No number appears more than once in a line. An example of an original grid and a solution are given below: Solve this problem by constraint logic programming. Each box of the matrix is considered as a domain, each box can take a value from 1 to 5. The constraints are simple since the sum of the boxes of a certain perimeter must be equal to a given number. L = [X11,X21,X31,X41,X51, X11 #= 5, X31*X41 #= 2, X32*X42 #= 4, X33 #= 4, X43*X53 #= 15, X34*X44 #= 15, X15*X25*X35*X45 #= 120, X21*X22 #= 15, X51*X52 #= 12, X12*X13*X14 #= 8, X23*X24 #= 2, X54*X55 #= 2, Skyscrapers are built on an NxN grid. Each skyscraper has a number of floors corresponding to its size (from 1 to N). In each column and each row each size appears exactly once. The numbers on the edges of the grid indicate how many skyscrapers can be seen when looking from that point towards the corresponding row or column. A skyscraper is visible if all the skyscrapers in front of it are smaller. Here is an example of a grid and its solution: For example, in the second line from the bottom looking from the left we can see 2 skyscrapers: the one of height 3 and the one of height 4. Here is a second problem of a grid (5×5) to solve: Solve this problem by constraint logic programming. You can use the predefined predicate fd_cardinality(+List, ?Count) where List is a list of constraints and Count is the number of List constraints satisfied. fd_cardinality([A5 #>A4 #/\A5 #>A3 #/\A5 #>A2 #/\A5 #>A1, A4 #>A3 #/\A4 #>A2 #/\A4 #>A1, A3 #> A2 #/\A3 #> A1, A2 #> A1],C). L = [Line1,Line2,Line3,Line4,Line5], Line1 = [X11,X12,X13,X14,X15], Line2 = [X21,X22,X23,X24,X25], Line3 = [X31,X32,X33,X34,X35], Line4 = [X41,X42,X43,X44,X45], Line5 = [X51,X52,X53,X54,X55], fd_all_different(Line3), fd_all_different(Line4), % The first skyscraper is still visible. So you have to remove 1 % of each number of visible skyscrapers. constraint(Line1,0), constraint(Line2,3), constraint(Line3,2),
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Elements of Geometry Because the polygon BECFD is greater than the circle LIM, the prism BCFKHG is greater than the cylinder LMNO, for they have the same altitude, but the prism has the greater base. But the pyramid ABECFD is the third part of the prism (15. 3. Sup.) BCFKHG, therefore it is greater than the third part of the cylinder LMNO. Now, the cone ABECFD is, by hypothesis, the third part of the cylinder LMNO, therefore the pyramid ABECFD is greater than the cone ABCD, and it is also less, because it is inscribed in the cone, which is impossible. Therefore, the cone ABCD is not less than the third part of the cylinder BFKG: And in the same manner, by circumscribing a polygon about the circle BCD, it may be shewn that the cone ABCD is not greater than the third part of the cylinder BFKG; therefore. it is equal to the third part of that cylinder. PROP. XIX. THEOR. If a hemisphere and a cone have equal bases and altitudes, a series of cylinders may be inscribed in the hemisphere, and another series may be described about the cone, having all the same altitudes with one another, and such that their sum shall differ from the sum of the hemisphere, and the cone, by a solid less than any given solid. Let ADB be a semicircle of which the centre is C, and let CD be at right angles to AB; let DB and DA be squares described on DC, draw CE, and let the figure thus constructed revolve about DC: then, the sector BCD, which is the half of the semicircle ADB, will describe a hemisphere having C for its centre (7 def. 3. Sup.), and the triangle CDE will describe a cone, having its vertex to C, and having for its base the circle (11. def. 3. Sup.) described by I)E, equal to that described by BC, which is the base of the hemisphere. Let W be any given solid. A series of cylinders may be inscribed in the hemisphere ADBĚ, and another described about the cone ECI, so that their sum shall differ from the sum of the hemisphere and the cone, by a solid less than the solid W. Upon the base of the hemisphere let a cylinder be constituted equal to W, and let its altitude be CX. Divide CD into such a number of equal parts, that each of them shall be less than CX; let these be CH, HG, GF, and FD. Through the points F, G, H, draw FN, GO, HP parallel to CB, meeting the circle in the points K, L and M; and the straight line CE in the points Q, R and S. From the points K, L, M draw Kf, Lg, Mh, perpendicular to GO, HP and CB; and from Q, R, and S, draw Qq, Rr, Ss, perpendicular to the same lines. It is evident, that the figure being thus constructed, if the whole revolve about CD, the rectangles Ff, Gg, Hh will describe cylinders (14. def. 3. Sup.) that will be circumscribed by the hemispheres BDA; and the rectangles DN, Fq, Gr, Hs, will also describe cylinders that will circumscribe the cone ICE. Now, it may be demonstrated, as was done of the prisms inscribed in a pyramid (13. 3. Sup.), that the sum of all the cylinders described within the hemisphere, is exceeded by the hemisphere by a solid less than the cylinder generated by the rectangle HB, that is, by a solid less than W, for the cylinder generated by HB is less than W. In the same manner, it may be demonstrated, that the sum of the cylinders circumscribing the cone ICE is greater than the cone by a solid less than the cylinder generated by the rectangle DN, that is, by a solid less than W. Therefore, since the sum of the cylinders inscribed in the hemisphere, together with a solid less than W, is equal to the hemisphere; and, since the sum of the cylinders described about the cone is equal to the cone together with a solid less than W; adding equals to equals, the sum of all these cylinders, together with a solid less than W, is equal to the sum of the hemisphere and the cone together with a solid less than W. Therefore, the difference between the whole of the cylinders and the sum of the hemisphere and the cone, is equal to the difference of two solids, which are each of them less than W; but this difference must also be less than W, therefore the difference between the two series of cylinders and the sum of the hemisphere and cone is less than the given solid W. PROP. XX. THEOR. The same things being supposed as in the last proposition, the sum of all the cylinders inscribed in the hemisphere, and described about the cone, is equal to a cylinder, having the same base and altitude with the hemisphere. Let the figure BCD be constructed as before, and supposed to revolve about CD; the cylinders inscribed in the hemisphere, that is, the cylinders described by the revolution of the rectangles Hh, Gg, Ff, together with those described about the cone, that is, the cylinders described by the revolution of the rectangles Hs, Gr, Fq, and DN are equal to the cylinder described by the revolution of the rectangle BD. Let L be the point in which GO meets the circle ABD, then, because CGL is a right angle if CL be joined, the circles described with the distances CG and GL are equal to the circle described with the distance CL (2. Cor. 6.1 Sup.) or GO; now, CG is equal to GR, because CD is equal to DE, and therefore also, the circles described with the distance GR and GL are together equal to the circle described with the distance GO, that is, the circles described by the revolution of GR and GL about the point G, are together equal to the circle described by the revolution of GO about the same point G; therefore also, the cylinders that stand upon the two first of these circles, having the common altitudes GH, are equal to the cylinder which stands on the remaining circle, and which has the same altitude GH. The cylinders described by the revolution of the rectangles Gg, and Gr are therefore equal to the cylinder described by the rectangle GP. And as the same may be shewn of all the rest, therefore the cylin ders described by the rectangles Hh, Gg, Ff, and by the rectangles Hs, Gr, Fq, DN, are together equal to the cylinder described by BD, that is, to the cylinder having the same base and altitude with the hemisphere. PROP. XXI. THEOR. Every sphere is two-thirds of the circumscribing cylinder. Let the figure be constructed as in the two last propositions, and if the hemisphere described by EDC be not equal to two-thirds of the cylinder described by BD, let it be greater by the solid W. Then, as the cone described by CDE is one-third of the cylinder (18. 3. Sup.) described by BD, the cone and the hemisphere together will exceed the cylinder by W. But that cylinder is equal to the sum of all the cylinders described by the rectangles Hh, Gg, Ff, Hs, Gr, Fq, DN (20. 3. Sup.); therefore the hemisphere and the cone added together exceed the sum of all these cylinders by the given solid W, which is absurd; for it has been shewn (19. 3. Sup.), that the hemisphere and the cone together differ from the sum of the cylinders by a solid less than W. The hemisphere is therefore equal to two-thirds of the cylinder described by the rectangle BD; and therefore the whole sphere is equal to two-thirds of the cylinder described by twice the rectangle BD, that is, to two-thirds of the circumscribing cylinder. END OF THE SUPPLEMENT TO THE ELEMENTS. PLANE TRIGONOMETRY. TRIGONOMETRY is the application of Arithmetic to Geometry: or, more precisely, it is the application of number to express the relations of the sides and angles of triangles to one another. It therefore necessarily supposes the elementary operations of arithmetic to be understood, and it borrows from that science several of the signs or characters which peculiarly belong to it. The elements of Plane Trigonometry, as laid down here, are divided into three sections: the first explains the principles; the second delivers the rules of calculation; the third contains the construction of trigonometrical tables, together with the investigation of some theorems, useful for extending trigonometry to the solution of the more difficult problems. SECTION I. LEMMA I. An angle at the centre of a circle is to four right angles as the arc on which it stands is to the whole circumference. Let ABC be an angle at the centre of the circle ACF, standing on the circumference AC: the angle ABC is to four right angles as the arc AC to the whole circumference ACF. Produce AB till it meet the circle in E, and draw DBF perpendicular to AE. Then, because ABC, ABD are two angles at the centre of the circle ACF, the angle ABC is to the angle ABD as the arc AC to the arc AD, (33. 6.); and therefore also, the angle ABC is to four times the angle ABD as the arc AC to four times the arc AD (4. 5.). But ABD is a right angle, and therefore four times the arc AD is equal to the whole circumference ACF; therefore the angle ABC is to four right angles as the arc AC to the whole circumference ACF. COR. Equal angles at the centres of different circles stand on arcs which have the same ratio to their circumferences. For, if the angle ABC, at the centre of the circles ACE, GHK, stand on the arcs AC, GH, AC is to the whole circumference of the circle ACE, as the angle ABC to four right angles; and the arc HG is to the whole circumference of the circle GHK in the same ratio. 1. IF two straight lines intersect one another in the centre of a circle, the arc of the circumference intercepted between them is called the Measure of the angle which they contain. Thus the arc AC is the measure of the angle ABC. 2. If the circumference of a circle be divided into 360 equal parts, each of these parts is called a Degree; and if a degree be divided into 60 equal parts, each of these is called a Minute; and if a minute be divided into 60 equal parts, each of them is called a Second, and so on. And as many degrees, minutes, seconds, &c. as are in any arc, so many degrees, minutes, seconds, &c. are said to be in the angle measured by that arc. COR. 1. Any arc is to the whole circumference of which it is a part, as the number of degrees, and parts of a degree contained in it is to the number 360. And any angle is to four right angles as the number of degrees and parts of a degree in the arc, which is the measure of that angle, is to 360. COR. 2. Hence also, the arcs which measure the same angle, whatever be the radii with which they are described, contain the same number of degrees, and parts of a degree. For the number of degrees and parts of a degree contained in each of these arcs has the same ratio to the number 360, that the angle which they measure has to four right angles (Cor. Lem. 1.). The degrees, minutes, seconds, &c. contained in any arc or angle, are usually written as in this example, 49°. 36′. 24′′. 42""; that is, 49 degrees, 36 minutes, 24 seconds, and 42 thirds. 3. Two angles, which are together equal to two right angles, or two arcs which are together equal to a semicircle, are called the Supplements of one another. 4. A straight line CD drawn through C, one of the extremities of the arc « ΠροηγούμενηΣυνέχεια »
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Vortices shed from a row of square bars Interactions between the wakes in a flow past a row of square bars are investigated by numerical simulations, the linear stability analysis and the bifurcation analysis. It is assumed that the row of square bars is placed across a uniform flow. Two-dimensional and incompressible flow field is also assumed. The flow is steady and symmetric along a streamwise centerline through the center of each square bar at low Reynolds numbers. However, it becomes unsteady and periodic in time at the Reynolds numbers larger than a critical value, and then the wakes behind the square bars become oscillatory. It is found by numerical simulations that vortices are shed synchronously from every couple of adjacent square bars in the same phase or in the anti-phase depending upon the distance between the bars. The synchronous shedding of vortices is clarified to occur due to an instability of the steady symmetric flow by the linear stability analysis. The bifurcation diagram of the flow is obtained and the critical Reynolds number of the instability is evaluated numerically. Translated title of the contribution Vortices shed from a row of square bars Original language Other Title of host publication 境界層遷移の解明と制御研究会講演論文集 Editors Takeshi Akinaga, Takao Kadoishi, Jiro Mizushima Pages 19-22 Number of pages 4 Publication status Published - Dec 2002 Publication series Name 航空宇宙技術研究所特別資料 Number NAL-SP-56 ISSN (Print) 1347-457X Dive into the research topics of 'Vortices shed from a row of square bars'. Together they form a unique fingerprint.
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Write Your Own map Method Now that you’ve seen how to write your own higher-order functions, let’s take a quick look at a more real-world example. Imagine for a moment that the List class doesn’t have its own map method, and you want to write your own. A good first step when creating functions is to accurately state the problem. Focusing only on a List[Int], you state: I want to write a map method that can be used to apply a function to each element in a List[Int] that it’s given, returning the transformed elements as a new list. Given that statement, you start to write the method signature. First, you know that you want to accept a function as a parameter, and that function should transform an Int into some generic type A, so you write: The syntax for using a generic type requires declaring that type symbol before the parameter list, so you add that: def map[A](f: (Int) => A) Next, you know that map should also accept a List[Int]: def map[A](f: (Int) => A, xs: List[Int]) Finally, you also know that map returns a transformed List that contains elements of the generic type A: def map[A](f: (Int) => A, xs: List[Int]): List[A] = ??? That takes care of the method signature. Now all you have to do is write the method body. A map method applies the function it’s given to every element in the list it’s given to produce a new, transformed list. One way to do this is with a for expression: for expressions often make code surprisingly simple, and for our purposes, that ends up being the entire method body. Putting it together with the method signature, you now have a standalone map method that works with a List[Int]: def map[A](f: (Int) => A, xs: List[Int]): List[A] = for (x <- xs) yield f(x) def map[A](f: (Int) => A, xs: List[Int]): List[A] = for x <- xs yield f(x) Make it generic As a bonus, notice that the for expression doesn’t do anything that depends on the type inside the List being Int. Therefore, you can replace Int in the type signature with the generic type parameter def map[A, B](f: (B) => A, xs: List[B]): List[A] = for (x <- xs) yield f(x) def map[A, B](f: (B) => A, xs: List[B]): List[A] = for x <- xs yield f(x) Now you have a map method that works with any List. These examples demonstrate that map works as desired: def double(i : Int): Int = i * 2 map(double, List(1, 2, 3)) // List(2, 4, 6) def strlen(s: String): Int = s.length map(strlen, List("a", "bb", "ccc")) // List(1, 2, 3) Now that you’ve seen how to write methods that accept functions as input parameters, let’s look at methods that return functions. Contributors to this page:
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efficiency of ball mill graph The graphs show that for a courser feed a three ball mixture (50 + 20 + 10 mm) will provide the much required largest amount of size class of interest. For a finer feed a binary mixture 20 mm + 10 mm will satisfy the objective function which is to get more of the 150 + 75 μm size class material.
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Whats the meaning of whole numbers? Whole numbers are a set of numbers including all positive integers and 0. Whole numbers are a part of real numbers that do not include fractions, decimals, or negative numbers. Counting numbers are also considered as whole numbers. What does whole numbers mean in math term? Whole Numbers The numbers that include natural numbers and zero. Not a fraction or decimal. What are whole numbers in short definition? A whole number is an exact number such as 1, 7, and 24, as opposed to a number with fractions or decimals. What’s another word for whole number? positive integers The whole numbers are also called the positive integers (or the nonnegative integers, if zero is included). How do you know if a number is a whole number? a more advanced example would be. You can multiply it by 10 and then do a “modulo” operation/divison with 10, and check if result of that two operations is zero. Result of that two operations will give you first digit after the decimal point. If result is equal to zero then the number is a whole number. How are whole numbers used in everyday life? Whole numbers like 0, 1, and 2 are the building blocks to understanding more complex number identifiers like real numbers, rational numbers, and irrational numbers. Rounding to the nearest whole number can also help you make calculations or do mental math more quickly. What is the difference between natural numbers and whole numbers? By definition, natural numbers are a part of the number system that contains all positive integers starting from the number 1 to infinity. Whereas, a whole number includes all positive numbers starting from the number 0 to infinity. The whole numbers are 121, 4, 0, 30. Which is the greatest whole number? So, 0, 1, 2, 3, 4…… are the whole numbers. We can clearly say that 1 is the smallest natural number and 0 is the smallest whole number. But there is no largest whole number because each number has its successor. Thus, there is no largest whole number. How many whole numbers are there in? My Standard Name Numbers Examples Whole Numbers { 0, 1, 2, 3, 4, } 0, 27,398, 2345 Counting Numbers { 1, 2, 3, 4, } 1, 18, 27, 2061 Integers { −4, −3, −2, −1, 0, 1, 2, 3, 4, } −15, 0, 27, 1102 What does whole numbers mean? Whole numbers are the same thing as counting numbers. Also known as natural numbers. They are used to count the number of physical objects. Example 101 students. What makes a whole number? A whole number is any number you can make by adding any number of 1s to 0 — including 0 itself. Some examples of whole numbers include 2, 5, 17, and 12,000. What are whole numbers examples? In mathematics, whole numbers are the basic counting numbers 0, 1, 2, 3, 4, 5, 6, … and so on. 17, 99, 267, 8107 and 999999999 are examples of whole numbers. What are some examples of a whole number? A whole number is any number that does not contain a fraction, decimal, or negative value. For example, 1, 25, and 365 are whole numbers.
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30000 BC - 500 BC - Chronology 30000 BC - 500 BC • Palaeolithic peoples in central Europe and France record numbers on bones. • Early geometric designs used. • A decimal number system is in use in Egypt. • Babylonian and Egyptian calendars in use. • The first symbols for numbers, simple straight lines, are used in Egypt. • The abacus is developed in the Middle East and in areas around the Mediterranean. • Hieroglyphic numerals in use in Egypt. (See this History Topic.) • Babylonians begin to use a sexagesimal number system for recording financial transactions. It is a place-value system without a zero place value. (See this History Topic.) • Harappans adopt a uniform decimal system of weights and measures. • The Moscow papyrus (also called the Golenishev papyrus) is written. It gives details of Egyptian geometry. (See this History Topic.) • Babylonians use multiplication tables. • The Babylonians solve linear and quadratic algebraic equations, compile tables of square and cube roots. They use Pythagoras's theorem and use mathematics to extend knowledge of astronomy. (See this History Topic.) • The Rhind papyrus (sometimes called the Ahmes papyrus) is written. It shows that Egyptian mathematics has developed many techniques to solve problems. Multiplication is based on repeated doubling, and division uses successive halving. (See this History Topic.) • About this date a decimal number system with no zero starts to be used in China. (See this History Topic.) • Apastamba writes the most interesting Indian Sulbasutra from a mathematical point of view. (See this History Topic.) • Thales brings Babylonian mathematical knowledge to Greece. He uses geometry to solve problems such as calculating the height of pyramids and the distance of ships from the shore. • Pythagoras of Samos moves to Croton in Italy and teaches mathematics, geometry, music, and reincarnation.
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Cite as Ruben Becker, Arnaud Casteigts, Pierluigi Crescenzi, Bojana Kodric, Malte Renken, Michael Raskin, and Viktor Zamaraev. Giant Components in Random Temporal Graphs. In Approximation, Randomization, and Combinatorial Optimization. Algorithms and Techniques (APPROX/RANDOM 2023). Leibniz International Proceedings in Informatics (LIPIcs), Volume 275, pp. 29:1-29:17, Schloss Dagstuhl – Leibniz-Zentrum für Informatik (2023) Copy BibTex To Clipboard author = {Becker, Ruben and Casteigts, Arnaud and Crescenzi, Pierluigi and Kodric, Bojana and Renken, Malte and Raskin, Michael and Zamaraev, Viktor}, title = {{Giant Components in Random Temporal Graphs}}, booktitle = {Approximation, Randomization, and Combinatorial Optimization. Algorithms and Techniques (APPROX/RANDOM 2023)}, pages = {29:1--29:17}, series = {Leibniz International Proceedings in Informatics (LIPIcs)}, ISBN = {978-3-95977-296-9}, ISSN = {1868-8969}, year = {2023}, volume = {275}, editor = {Megow, Nicole and Smith, Adam}, publisher = {Schloss Dagstuhl -- Leibniz-Zentrum f{\"u}r Informatik}, address = {Dagstuhl, Germany}, URL = {https://drops.dagstuhl.de/entities/document/10.4230/LIPIcs.APPROX/RANDOM.2023.29}, URN = {urn:nbn:de:0030-drops-188542}, doi = {10.4230/LIPIcs.APPROX/RANDOM.2023.29}, annote = {Keywords: random temporal graph, Erd\H{o}s–R\'{e}nyi random graph, sharp threshold, temporal connectivity, temporal connected component, edge-ordered graph}
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Math 116 Course Objectives: This is a fundamental course in numerical analysis, with a focus on numerical linear algebra and optimzation methods, which are the building blocks for a variety of algorithms used in data science, machine learning, signal and image processing, and evolutionary dynamics. The primary goal is for students to gain a complete understanding of how to evaluate the efficacy of (general) numerical methods, specifically their accuracy, stability, and convergence properties. This is not a survey course. Students will write simple computer programs in MATLAB or python to gain understanding of how methods work, in particular to compare advantages and disadvantages of various algorithms. This course will not emphasize programming or the use of numerical software. Numerical analysis is a fundamental topic in applied mathematics. Many practical problems that scientists try to solve are based on mathematical models, but few can be solved analytically, either due to their complexity, large number of variables, or lack of information. Computational algorithms are therefore needed for approximating these solutions. It is critically important to maintain the important mathematical properties of the underlying system when developing these computational algorithms. Numerical analysis is about developing good computational techniques for broad based problems and analyzing their properties. Numerical analysts seek to demonstrate when computational algorithms are trustworthy so that domain scientists can be confident in the results of their experiments. Numerical analysts also answer the question, "what assumptions of the underlying problem are necessary for this computational method to succeed?"
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Binning data by quantiles? Beware of rounded data In my article about how to create a quantile plot, I chose not to discuss a theoretical issue that occasionally occurs. The issue is that for discrete data (which includes rounded values), it might be impossible to use quantile values to split the data into k groups where each group has approximately the same number of elements. To put it bluntly, the algorithm that I proposed for creating a quantile bin plot can fail for large values of k. The problem can occur for data that are rounded to the nearest unit, such as age, height, weight, and blood pressure. The algorithm assumes that the data quantiles are (mostly) unique and that they divide the empirical cumulative distribution function (ECDF) of N observations into k groups that each have approximately N/k observations. However, if there are more than N/k repeated values, the repeated value can occupy more than one quantile value. In fact, this will always happen if a particular value is repeated more than 2N/k times. For example, suppose that you want to divide the diastolic blood pressure data in the Sashelp.Heart data set into 10 groups that each have approximately 10% of the data. The following statements draw the ECDF for the diastolic measurements and mark the deciles: proc univariate data=Sashelp.Heart; var diastolic; cdfplot diastolic / vref=(10 to 90 by 10); ods select CDFPlot Quantiles; The value 80 appears in 711 of the 5209 observations, which is about 14% of the data. Notice that the value 80 intersects two reference lines, one at 30 and the other at 40. This means that 80 is both the 30th percentile and the 40th percentile. Assuming that you do not want to split the value 80 across two bins, the deciles of the data produce at most nine bins. Another way to see this graphically is to use the RANK procedure to try to group the data into 10 groups, as described in the article "Grouping observations based on quantiles." The following statements create a new variable called Group, which for continuous data would have the values 0–9. However, for this rounded data only nine groups exist: proc sort data=Sashelp.Heart out=Heart; by Diastolic; proc rank data=Heart out=Heart groups=10 ties=high; var Diastolic; /* variable on which to group */ ranks Group; /* name of variable to contain groups 0,1,...,k-1 */ proc sgplot data=Heart; scatter x=Group y=Diastolic / group=Group; refline 80 / axis=y; xaxis values=(0 to 9); As shown by the legend on the scatter plot, only nine groups were created, instead of the 10 that were asked for. Group 3 was not created. I've described a problem, but what can you do about it? Not a lot. The problem is in the data. If you are flexible about the number of groups, you can try decreasing the value of k. Because each group contains approximately N/k observations, decreasing k might circumvent the problem. Another possibility is to jitter the data by adding a small amount of random noise. That will eliminate the problem of repeated values. So I'll leave you with this warning: beware of using quantiles to bin rounded data into groups. Although the technique works great when almost all of the data values are distinct, you can run into problems if you ask for many bins and your data contain many repeated values. 2 Comments Leave A Reply
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division Archives - Page 2 of 4 - EducationUnboxed.com - Free Help for Homeschool This is part three in the series of videos about long division. This free math tutoring video has divisors and quotients that are 2-digit numbers and shows how you can use drawings to aid in the conceptual understanding of these types of problems. This method works especially well for visual, right-brained, or global learners. For more practice, click here to download a worksheet. Long Division – Part Two – Free Math Tutoring Video This is the second free math tutoring video in the long division series. These problems have divisors and quotients that are between 10 and 19. Doing multiplication problems with Cuisenaire Rods where they multiply teen numbers would be an excellent intro to this topic. This method enables children to truly understand what the long division algorithm is all about instead of simply memorizing the steps. It is great for visual and kinesthetic learners! For more practice, click here to download a worksheet. Long Division – Part One – Free Math Tutoring Video This is the first of several free math tutoring videos that show how to do long division using Cuisenaire Rods to aid in conceptual understanding. Most people have no idea WHY the long division algorithm works. They’ve just memorized the steps without understanding. This may work for many people, but there are a whole lot of children in school today who have sequencing issues, who struggle with memorizing the steps to math formulas that have no meaning to them, who simply want to understand WHY all this stuff works. Those children will be better served if their teachers understand why and can teach from a conceptual framework instead of a formula-driven approach. Even for those students who are willing to just “memorize the formula,” this method will be beneficial because it builds a clearer understanding of the decimal system and the distributive property which will give them confidence that they really do understand and are “good at” math. I discovered this way of teaching long division through Crewton Ramone’s House of Math. I’ve changed his method very slightly to what makes more sense to me, but it is essentially the same idea. This method is ingenious. I don’t know why more people don’t know about it! For more practice, click here to download a worksheet.
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• A Sharp Theory of Hardy Spaces on Ahlfors-Regular Quasi-Metric Spaces. With M.Mitrea, Springer Lecture Notes in Mathematics, Vol.2142, (2015), 486 pages. Introduction Springerlink • Borel regularity is equivalent to Lusin's theorem and the existence of Borel representatives. With P.Górka and A.Słabuszewski. (submitted) arXiv • Compact embeddings of Sobolev, Besov, and Triebel-Lizorkin spaces. With P.Górka and A.Słabuszewski. (submitted) arXiv • Optimal Embeddings for Triebel-Lizorkin and Besov Spaces on Quasi-Metric Measure Spaces. With D.Yang and W.Yuan, Math. Z. 307, 50 (2024). arXiv • A simple proof of reflexivity and separability of N^{1,p} Sobolev spaces. With P.Hajłasz and Lukáš Malý. Ann. Fenn. Math. 48 (2023), no. 1, 255-275. arXiv • Pointwise Characterization of Besov and Triebel–Lizorkin Spaces on Spaces of Homogeneous Type. With F.Wang, D.Yang and W.Yuan. Studia Math. 268 (2023), no. 2, 121–166. PDF • A Measure Characterization of Embedding and Extension Domains for Sobolev, Triebel-Lizorkin, and Besov Spaces on Spaces of Homogeneous Type. With D.Yang and W.Yuan. J. Funct. Anal. 283 (2022), no. 12, Paper No. 109687, 71 pp. arXiv • The game of cycles. With M.Averett, B.Gaines, C.Jackson, M.L.Karker, M.A.Marciniak, F.E.Su, and S.Walker. Amer. Math. Monthly 128 (2021), no. 10, 868-887. arXiv • Sobolev embedding for M^{1,p} spaces is equivalent to a lower bound of the measure. With P.Górka and P.Hajłasz. J. Funct. Anal. 279 (2020), no. 7, 108628 arXiv • A note on metric-measure spaces supporting Poincaré inequalities. With P.Hajłasz. Atti Accad. Naz. Lincei Rend. Lincei Mat. Appl. 31 (2020), no.1, 15-23. arXiv • Characterizing Lusin's Theorem and the Density of Continuous Functions in Lebesgue Spaces via the Regularity of the Measure. With M.Mitrea and B.Schmutzler. (preprint) • Whitney-type extensions with control of the modulus of continuity in geometrically doubling quasi-metric spaces. With I.Mitrea and M.Mitrea. Commun. Pure Appl. Anal., 12 (2013), No. 1, pp. 59-88. • Sharp geometric maximum principles for semi-elliptic operators with singular drift. With D.Brigham, V.Maz'ya, M.Mitrea and E.Ziadé, Math. Res. Lett., Vol.18 (2011), No. 04, pp. 613-620. PDF • On the regularity of domains satisfying a uniform hour-glass condition and a sharp version of the Hopf-Oleinik Boundary Point Principle. With D.Brigham, V.Maz'ya, M.Mitrea and E.Ziadé, Journal of Mathematical Sciences, Vol. 176 (2011), No. 3, pp. 281-360. • Topics in Harmonic Analysis and Partial Differential Equations: Extension Theorems and Geometric Maximum Principles. Masters Thesis. (2011). PDF
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Automata Theory Questions for Campus Interviews - Sanfoundry Automata Theory Questions and Answers – Pumping Lemma for Context Free Language This set of Automata Theory Questions and Answers for Campus interviews focuses on “Pumping Lemma for Context Free Language”. 1. Which of the following is called Bar-Hillel lemma? a) Pumping lemma for regular language b) Pumping lemma for context free languages c) Pumping lemma for context sensitive languages d) None of the mentioned View Answer Answer: b Explanation: In automata theory, the pumping lemma for context free languages, also kmown as the Bar-Hillel lemma, represents a property of all context free languages. 2. Which of the expressions correctly is an requirement of the pumping lemma for the context free languages? a) uv^nwx^ny b) uv^nw^nx^ny c) uv^2nwx^2ny d) All of the mentioned View Answer Answer: b Explanation: Let L be a CFL. Then there is an integer n so that for any u that belong to language L satisfying |t| >=n, there are strings u, v, w, x, y and z satisfying |vwx|<=n For any m>=0, uv^nwx^ny ∈ L 3.Let L be a CFL. Then there is an integer n so that for any u that belong to language L satisfying |t|>=n, there are strings u, v, w, x, y and z satisfying Let p be the number of variables in CNF form of the context free grammar. The value of n in terms of p is a) 2p b) 2p c) 2p+1 d) p^2 View Answer Answer: c Explanation: This inequation has been derived from derivation tree for t which must have height at least p+2(It has more than 2^p leaf nodes, and therefore its height is >p+1). 4. Which of the following gives a positive result to the pumping lemma restrictions and requirements? a) {a^ib^ic^i|i>=0} b) {0^i1^i|i>=0} c) {ss|s∈{a,b}*} d) None of the mentioned View Answer Answer: b Explanation: A positive result to the pumping lemma shows that the language is a CFL and ist contradiction or negative result shows that the given language is not a Context Free language. 5. Using pumping lemma, which of the following cannot be proved as ‘not a CFL’? a) {a^ib^ic^i|i>=0} b) {ss|s∈{a,b}*} c) The set legal C programs d) None of the mentioned View Answer Answer: d Explanation: There are few rules in C that are context dependent. For example, declaration of a variable before it can be used. 6. State true or false: Statement: We cannot use Ogden’s lemma when pumping lemma fails. a) true b) false View Answer Answer: b Explanation: Although the pumping lemma provides some information about v and x that are pumped, it says little about the location of these substrings in the string t. It can be used whenever the pumping lemma fails. Example: {a^pb^qc^rd^s|p=0 or q=r=s}, etc. 7. Which of the following cannot be filled in the blank below? Statement: There are CFLs L1 nad L2 so that ___________is not a CFL. a) L1∩L2 b) L1′ c) L1* d) None of the mentioned View Answer Answer: c Explanation: A set of context free language is closed under the following operations: a) Union b) Concatenation c) Kleene 8. The pumping lemma is often used to prove that a language is: a) Context free b) Not context free c) Regular d) None of the mentioned View Answer Answer: b Explanation: The pumping lemma is often used to prove that a given language L is non-context-free, by showing that arbitrarily long strings s are in L that cannot be “pumped” without producing strings outside L. 9. What is the pumping length of string of length x? a) x+1 b) x c) x-1 d) x2 View Answer Answer: a Explanation: There exists a property of all strings in the language that are of length p, where p is the constant-called the pumping length .For a finite language L, p is equal to the maximum string length in L plus 1. 10. Which of the following does not obey pumping lemma for context free languages ? a) Finite languages b) Context free languages c) Unrestricted languages d) None of the mentioned View Answer Answer: c Explanation: Finite languages (which are regular hence context free ) obey pumping lemma where as unrestricted languages like recursive languages do not obey pumping lemma for context free languages. Sanfoundry Global Education & Learning Series – Automata Theory. To practice all areas of Automata Theory for Campus Interviews, here is complete set of 1000+ Multiple Choice Questions and Answers.
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Cite as Mikołaj Bojańczyk, Lê Thành Dũng (Tito) Nguyễn, and Rafał Stefański. Function Spaces for Orbit-Finite Sets. In 51st International Colloquium on Automata, Languages, and Programming (ICALP 2024). Leibniz International Proceedings in Informatics (LIPIcs), Volume 297, pp. 130:1-130:20, Schloss Dagstuhl – Leibniz-Zentrum für Informatik (2024) Copy BibTex To Clipboard author = {Boja\'{n}czyk, Miko{\l}aj and Nguy\~{ê}n, L\^{e} Th\`{a}nh D\~{u}ng (Tito) and Stefa\'{n}ski, Rafa{\l}}, title = {{Function Spaces for Orbit-Finite Sets}}, booktitle = {51st International Colloquium on Automata, Languages, and Programming (ICALP 2024)}, pages = {130:1--130:20}, series = {Leibniz International Proceedings in Informatics (LIPIcs)}, ISBN = {978-3-95977-322-5}, ISSN = {1868-8969}, year = {2024}, volume = {297}, editor = {Bringmann, Karl and Grohe, Martin and Puppis, Gabriele and Svensson, Ola}, publisher = {Schloss Dagstuhl -- Leibniz-Zentrum f{\"u}r Informatik}, address = {Dagstuhl, Germany}, URL = {https://drops.dagstuhl.de/entities/document/10.4230/LIPIcs.ICALP.2024.130}, URN = {urn:nbn:de:0030-drops-202730}, doi = {10.4230/LIPIcs.ICALP.2024.130}, annote = {Keywords: Orbit-finite sets, automata, linear types, game semantics}
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mathematics Archives • Turton School Induction to KS4 A demonstration of Sparx Maths and Corbett Maths with Mr. Hodgson. Induction to KS4 A demonstration of Sparx Maths and Corbett Maths with Mr. Hodgson. Please find additional guides for other Year 10 subjects here… English Science Faith & Ethics Year 11 Science & Maths Exams (Mocks /OCT) Please find attached a copy of the exam timetable for the upcoming Year 11 Science and Maths exams. All year 11 students have been issued with this timetable and have a copy on their school email account. Your son / daughter’s Maths and Science teachers will have discussed these exams with them and all Read More … Year 11 Revision Outline Maths Higher (Mathematics) Year 11 Revision Outline for Maths Higher (Mathematics) Year 11 Revision Outline for Maths Foundation (Mathematics) Year 11 Revision Outline for Maths Foundation (Mathematics) Maths Learning Walk Maths Learning WalkOn 5th January 2016 fellow Governor, Dave Miller and I met with Deputy Head Teacher, Mrs Bach and Ms Gorse the Head Teacher. Initially we observed their book scrutiny of a selection of pupils’ work books from various classes and years. Their very thorough scrutiny was to ensure consistent teaching and marking was Read More …
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Suggestion: support graphing families of functions Topic: Suggestion: support graphing families of functions A friend just showed me this package last night to illustrate some functions we were looking at. Quite liked it. But when I asked her if she could plot a family of functions, no can do. Pretty important for looking at trends. To clarify, imagine the very simple function: y = x^n it would be nice to ask Graph to plot this for n=0 to 5 step 1 say. Thus a family of functions varying n. The same applies to more complex functions indeed is where it becomes truly useful, to look at how a given coefficient alters the plot. As it was she had to create a separate line manually for each value we wanted to look at, time consuming. Re: Suggestion: support graphing families of functions Animations can do something a little similar. But instead of adding more functions it will change the existing one. That is you can add the function f(x)=x^n and animate n from 0 to 5. A more advanced possibility is to use the scripting engine, which requires you to install Python 3.2 32 bit. When you have Python installed you can press F11 inside Graph to get a Python interpreter and enter the following code to add the family of functions: for n in range(6): F=Graph.TStdFunc("x^%d" % n) I will consider creating an easier way to do this in a future version.
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Roots third order quadratic roots third order quadratic Related topics: math power 8 algebra program algebra answer solver products and quotients of ration intermediate algebra answers Rational Expression Solver positive and negative integer worksheets inequality in excel radical calculators quadratic equation matlab gcse mcqs on maths physics chemistry free online algebra test free algebrator Least Common Denominator In Algebra Author Message Calnicta Posted: Wednesday 07th of Dec 07:13 I have trouble with roots third order quadratic. I tried a lot to get somebody who can assist me with this. I also searched for a coach to teach me and crack my problems on quadratic formula, simplifying expressions and radicals. Though I found a few who could perhaps crack my problem, I recognized that I cannot afford them. I do not have much time too. My exam is coming up in a little while. I am worried . Can someone assist me with this situation? I would really value any help or any information. From: UK Back to top IlbendF Posted: Wednesday 07th of Dec 13:12 Algebrator is one of the most powerful resources that can offer help to a person like you. When I was a beginner, I took support from Algebrator. Algebrator offers all the basics of Remedial Algebra. Rather than using the Algebrator as a step-by-step guide to solve all your homework assignments, you can use it as a coach that can offer the basics of radicals, powers and trinomials. Once you assimilate the principles, you can go ahead and work out any tough assignments on Algebra 1 in no time . Back to top Ashe Posted: Thursday 08th of Dec 19:58 Algebrator is one beneficial tool. I don’t have much interest in math and have found it to be complicated all my life. Yet one cannot always leave math because it sometimes becomes a compulsory part of one’s course work. My friend is a math wiz and I found this software in his palmtop . It was only then I understood why he finds this subject to be so simple . Back to top vromolx Posted: Saturday 10th of Dec 08:19 I'm so happy I got these answers so fast, I can't wait to buy Algebrator. Can you tell me one more thing, where could I find this program? I'm not so good at searching for things like this, so it would be good if you could give me a link . Thanks a lot! From: Alken - Belgium Back to top Vild Posted: Saturday 10th of Dec 14:31 Here is the link https://softmath.com/comparison-algebra-homework.html Back to top Bet Posted: Saturday 10th of Dec 16:45 Algebrator is the program that I have used through several math classes - College Algebra, College Algebra and Algebra 1. It is a really a great piece of math software. I remember of going through difficulties with graphing, trinomials and radical expressions. I would simply type in a problem from the workbook , click on Solve – and step by step solution to my math homework. I highly recommend the program. From: kµlt øƒ Ø™ Back to top
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Graphical Representation of Motion along a Straight Line Plotting the distance/displacement or speed/velocity on a graph helps us visually understand certain things about time and position. 1. The distance – time graph for uniform motion The following Table shows the distance walked by Surya at different times. A graph is drawn by taking time along X-axis and distance along Y-axis. The graph is known as distance – time graph. When we look at the distance – time graph of Surya’s walk, we notice certain things. First, it is a straight line. We also notice that Surya covers equal distances in equal intervals of time. We can therefore conclude that Surya walked at a constant speed. Can you find the speed at which Surya walked, from the graph? Yes, you can. The parameter is referred to as the slope of the line. Speed at which Surya walked = distance covered / time taken = BC/AC (From the graph) = slope of the straight line = 500 / 5 = 100 ms^-1 Steeper the slope (in other words the larger value) the greater is the speed. Let us take a look at the distance–time graphs of three different people – Surya walking, Monica cycling and Hari going in a car, along the same path. We know that cycling can be faster than walking and a car can go faster than a cycle. The distance – time graph of the three would be as given in the following graph. The slope of the line on the distance – time graph becomes steeper and steeper as the speed increases. 2. The distance time graph for non uniform motion We can also plot the distance – time graph for accelerated motion (non uniform motion). Table given below shows the distance travelled by a car in a time interval of two second. Note that the graph is not a straight line as we got in the case of uniform motion. is nature of the graph shows non – linear variation of the distance travelled by the car with time. us, the graph represents motion with non uniform speed. 3. Velocity – Time graph The variation in velocity of an object with time can be represented by velocity – time graph. In the graph, time is represented along the X – axis and the velocity is represented along the Y – axis. If the object moves at uniform velocity, a straight line parallel to X-axis is obtained. is Graph shows the velocity – time graph for a car moving with uniform velocity of 40 km/hour. We know that the product of velocity and time gives displacement of an object moving with uniform velocity. The area under the velocity – time graph is equal to the magnitude of the displacement. So the distance (displacement) S covered by the car in a time interval of t can be expressed as S = AC × CD S = Area of the rectangle ABCD (shaded portion in the graph) We can also study about uniformly accelerated motion by plotting its velocity – time graph. Consider a car being driven along a straight road for testing its engine. Suppose a person sitting next to the driver records its velocity for every 5 seconds from the speedometer of the car. The velocity of the car in ms-1 at different instants of time is shown in the Table below. In this case, the velocity – time graph for the motion of the car is shown in graph (straight line). The nature of the graph shows that the velocity changes by equal amounts in equal intervals of time. Thus, for all uniformly accelerated motion, the velocity – time graph is a straight line. One can also determine the distance moved by the car from its velocity – time graph. The area under the velocity – time graph gives the distance (magnitude of displacement) moved by the car in a given interval of time. Since the magnitude of the velocity of the car is changing due to acceleration, the distance S travelled by the car will be given by the area ABCDE under the velocity – time graph. That is S = area ABCDE = area of the rectangle ABCD + area of the triangle ADE S = (AB × BC) + ½ (AD × DE) The area ABCDE can also be calculated by considering the shape as trapezium. Area of the quadrangle ABCDE can also be calculated by calculating the area of trapezium ABCDE. It means S = area of trapezium ABCDE = ½ × sum of length of parallel sides × distance between parallel sides S = ½ × (AB + CE) × BC In the case of non uniformly accelerated motion, distance – time graph, velocity – time graphs can have any shape as shown in Figure below:
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Ch 7-8 Review Questions Review Questions Ch 7 & 8 – Work and Energy Concept Questions c1. Sign of work done In the following scenarios, a force pushes on an object. Determine if the work done by the force is positive, negative, or zero. a. A box is pushed toward the left, and slides along the floor at a constant speed toward the left. b. The same box as in (a). Consider the work done by gravity and work done by the normal force. c. A car turns a corner on a level road at constant speed. Consider the work done by static friction, which causes the centripetal acceleration. d. A crane lowers a mass at constant speed. Consider the work by the tension in the cable. e. Force \(\vec F\) and displacement \(\Delta \vec r\), as shown below. a. Work done by pushing force is positive since force and displacement are in the same (leftward) direction. b. Work done by gravity and the normal force are zero, because these forces are perpendicular to the motion. c. Work done by static friction is zero because \(\vec f\) is always perpendicular to \(\vec v\). d. Work done by the tension is negative because it points up while the motion is down. e. Work done is postive because \(\vec F \cdot \Delta \vec r > 0\). c2. Energy Transformations For the following scenarios, describe the energy transformation taking place. Use \(K\), \(U\), and \(E_{th}\) to refer to kinetic, potential and thermal energy, respectively. a. A diver falls downward, before hitting the water. b. The diver hits the water and rapidly slows. c. A skier is sliding down a slope at a constant speed. d. A box slides up a gentle but slightly rough incline until stopping at the top. a. \(U \rightarrow K\) b. \(K \rightarrow E_{th}\) c. \(U \rightarrow E_{th}\) d. \(K \rightarrow U\) and \(E_{th}\) a. As the diver falls, they pick up speed. Elevation loss means \(U\) decreases. Speed increase means \(K\) increases. b. The speed decrease means \(K\) decreases. There is no substantial elevation change, so the energy has been lost to thermal energy in the water. c. The skier’s speed is constant, so \(K\) does not change. There is a loss of elevation, meaning \(U\) is lost. If the skier does not speed up, it must be a friction force that does work. Meaning thermal energy is generated. d. Losing speed means \(K\) is lost. Gaining elevation means \(U\) increases. The “rough” surface implies that friction does work, so thermal energy is also created. c3. Work and kinetic energy Object \(A\) has mass \(M\), and \(B\) has mass \(2M\). Both are at rest. Compare their kinetic energies (\(K_A\) and \(K_B\)) after a. They are each pushed with the same force for the same distance \(d\). b. They are each pushed with the same force for the same time \(t\). a. \(K_A = K_B\) b. \(K_A \gt K_B\) In each case, \(K\) depends on the work done, which is force times distance. a. Work done is the same, so \(\Delta K\) is the same. b. Over time \(t\), \(A\) will accelerate more and move a greater distance. So work done on \(A\) is greater. c4. Height vs speed A roller coaster rolls from rest down a frictionless track, reaching speed \(v\) at the bottom. If you want the car to go twice as fast at the bottom, by what factor must you increase the height of the track? The track must be \(4\) times higher. Energy conservation gives \(mgh = \frac{1}{2} m v^2\), so \(h = \frac{v^2}{2g}\). So if \(v\) is doubled, the new height will be \(4h\). c5. Spring energy Rank the elastic potential enegies stored in the following springs. \(U_1 = \frac{1}{2} k d^2\) \(U_2 = \frac{1}{2} k d^2 = U_1\) \(U_3 = \frac{1}{2} (2k) d^2 = 2 U_1\) \(U_4 = \frac{1}{2} k (2d)^2 = 4 U_1\) \(U_4 > U_3 > U_2 = U_1\) c6. Spring gun muzzle velocity A spring-loaded gun shoots out a plastic ball at speed \(v_0\). The spring is then compressed twice the distance it was on the first shot. a. By what factor is the spring’s potential energy increased? b. By what factor is the ball’s launch speed increased? a. \(U = \frac{1}{2} k x^2\), so double \(x\) means four times more energy. b. \(U \rightarrow \frac{1}{2} m v^2\), so four times more energy means double \(v\). c7. Energy diagram Below are a set of axes on which you are going to draw a potential-energy curve. By doing experiments, you find the following information: • A particle with energy \(E_1\) oscillates between positions D and E. • A particle with energy \(E_2\) oscillates between positions C and F. • A particle with energy \(E_3\) oscillates between positions B and G. • A particle with energy \(E_4\) enters from the right, turns around at A, then never returns. Draw a potential-energy curve that is consistent with this information. c8. Features of energy diagram The graph shows the potential energy curve of a particle moving along the \(x\)-axis under the influence of a conservative force. a. In which intervals of \(x\) is the force on the particle to the right? b. In which intervals of \(x\) is the force on the particle to the left? c. At what values of \(x\) is the magnitude of the force maximum? d. What values of \(x\) are positions of stable equilibrium? e. What values of \(x\) are positions of unstable equilibrium? f. If the particle is released from rest at \(x=0\) m, will it reach \(x=10\) m? a. \(F_x= - (dU)/dx\) (negative of slope). So \(\vec F_x\) is to the right from \(0\) m to \(2\) m, and from \(5\) m to \(8\) m. b. \(\vec F_x\) is to the left from \(2\) m to \(5\) m, and from \(8\) m to \(10\) m. c. Slope is steepest at \(x=0\) m, \(x=3.5\) m, \(x=6.5\) m, and \(x=10\) m. d. At the minima, \(x=2\) m and \(x=8\) m. e. At the maxima, \(x=5\) m. f. No. The particle will move right until it turns around at about \(x=4\) m. c9. Compare projectiles Two projectiles (with the same mass) are launched over flat ground with the same intial speed. Projectile \(A\) is launched at angle \(10\)°, and projectile \(B\) at angle \(50\)° (above horizontal). You may ignore air resistance. a. Which projectile has greater speed at its peak height? b. Which projectile has greater speed upon hitting the ground? Note the equal initial speed means that they have the same total energy at all times throughout their flights: \(K+U =\) constant. a. \(A\) is moving faster. \(K_A \gt K_B\), since \(U_A \lt U_B\) at their peak height. b. Same speed. \(K_A = K_B\), since \(U_A=U_B=0\) at the ground. Long Answer Questions 1. Work done by given force Consider a particle on which several forces act, one of which is known to be constant in time: \(\vec{F}_1 = 3 \hat{\imath} + 4 \hat{\jmath}\) N As a result, the particle moves along a straight path from a Cartesian coordinate of (\(-1\) m, \(2\) m) to (\(5\) m, \(6\) m). What is the work done by \(\vec F_1\) ? Since the force is constant and the path is straight, \(W = \vec F \cdot \Delta \vec s\). \(W= (3 \hat{\imath} + 4 \hat{\jmath}) \cdot (6\hat{\imath} + 4 \hat{\jmath}) = 18+16 =34\) N·m 2. Accelerated electron An electron in a television tube is accelerated uniformly from rest to a speed of \(8.4 \times 10^7\) m/s over a distance of \(2.5\) cm. What is the power delivered to the electron at the instant that its displacement is \(1.0\) cm? The electron has a mass of \(9.11 \times 10^{-31}\) kg. Work done is the gain in kinetic energy: \(W=Fd= \Delta K =\frac{1}{2} m v^2\). So with \(d=2.5\) cm, you can find the force on the electron. From this, acceleration is \(a = F/m = 1.41 \times 10^{17}\) m/s² Now with \(d= 1.0\) cm, find velocity with \(v^2 = 2 a \Delta s\). \(v = 5.31 \times 10^{7}\) m/s Then \(P = F \cdot v\). \(P = 6.83\) μW 3. Work done by different forces A \(40\) kg crate is pushed uphill at constant velocity a distance \(8.0\) m along a \(30\)° incline by the horizontal force \(\vec F\). The coefficient of kinetic friction between the crate and the incline is \(\mu_k\) = \(0.40\). Calculate the work done by a. the applied force b. the normal force c. the frictional force d. the gravitational force e. the net force a. \(3.46\) kJ b. \(0\) c. \(–1.89\) kJ d. \(–1.57\) kJ e. \(0\) Consider the free-body diagram: Note that the normal force is perpendicular to displacement, and therefore does zero work. \(\vec{F}=m\vec{a}\) with the axes shown gives the following: \(\Sigma F_x = 0 = F - N \sin \theta - \mu N \cos \theta\) \(\Sigma F_y = 0 = N \cos \theta - \mu N \sin \theta - mg\) where \(\theta = 30^{\circ}\). The second equation yeilds \(N = \frac{mg}{\cos \theta - \mu \sin \theta} = 588.57\) N. With this, the first equation yeilds \(F = 498.17\) N. With these you can find work done with \(W = \vec F \cdot \Delta \vec s\). Since \(\vec{a} = 0\), \(\vec{F}_{\textup{net}} = 0\). So the net force does zero work. 4. Kinetic Energy after work is done A 10-kg box slides along a track at 8 m/s. As the box slides up a ramp, it gains \(1.0\) m of elevation, while friction does \(-100\) J of work on it. How fast is it moving afterward? The change in energy \(\Delta (K+U)\) equals work done by friction (which is negative because energy is lost). \(K_f + U_f - (K_i +U_i) = W\) \(K_f = \frac{1}{2} m v_f^2\) \(U_f = mgh\) \(K_i = \frac{1}{2} m v_i^2\) \(U_i = 0\) Plug these in and solve for \(v_f\). \(v_f^2 = v_i^2 - 2gh + 2\frac{W}{m} = 8^2 - 2(9.8)(1.0) - 2(100)/10 = 24.4\) m²/s² \(v_f = 4.9\) m/s Last modified: Tue October 22 2024, 09:25 PM.
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Femtometer/Second Squared Converter | Kody Tools Conversion Description 1 Femtometer/Second Squared = 1e-15 Meter/Second Squared 1 Femtometer/Second Squared in Meter/Second Squared is equal to 1e-15 1 Femtometer/Second Squared = 1000 Attometer/Second Squared 1 Femtometer/Second Squared in Attometer/Second Squared is equal to 1000 1 Femtometer/Second Squared = 1e-13 Centimeter/Second Squared 1 Femtometer/Second Squared in Centimeter/Second Squared is equal to 1e-13 1 Femtometer/Second Squared = 1e-14 Decimeter/Second Squared 1 Femtometer/Second Squared in Decimeter/Second Squared is equal to 1e-14 1 Femtometer/Second Squared = 1e-16 Dekameter/Second Squared 1 Femtometer/Second Squared in Dekameter/Second Squared is equal to 1e-16 1 Femtometer/Second Squared = 1e-17 Hectometer/Second Squared 1 Femtometer/Second Squared in Hectometer/Second Squared is equal to 1e-17 1 Femtometer/Second Squared = 1e-18 Kilometer/Second Squared 1 Femtometer/Second Squared in Kilometer/Second Squared is equal to 1e-18 1 Femtometer/Second Squared = 1e-9 Micrometer/Second Squared 1 Femtometer/Second Squared in Micrometer/Second Squared is equal to 1e-9 1 Femtometer/Second Squared = 1e-12 Millimeter/Second Squared 1 Femtometer/Second Squared in Millimeter/Second Squared is equal to 1e-12 1 Femtometer/Second Squared = 0.000001 Nanometer/Second Squared 1 Femtometer/Second Squared in Nanometer/Second Squared is equal to 0.000001 1 Femtometer/Second Squared = 0.001 Picometer/Second Squared 1 Femtometer/Second Squared in Picometer/Second Squared is equal to 0.001 1 Femtometer/Second Squared = 1.296e-8 Meter/Hour Squared 1 Femtometer/Second Squared in Meter/Hour Squared is equal to 1.296e-8 1 Femtometer/Second Squared = 0.00001296 Millimeter/Hour Squared 1 Femtometer/Second Squared in Millimeter/Hour Squared is equal to 0.00001296 1 Femtometer/Second Squared = 0.000001296 Centimeter/Hour Squared 1 Femtometer/Second Squared in Centimeter/Hour Squared is equal to 0.000001296 1 Femtometer/Second Squared = 1.296e-11 Kilometer/Hour Squared 1 Femtometer/Second Squared in Kilometer/Hour Squared is equal to 1.296e-11 1 Femtometer/Second Squared = 3.6e-12 Meter/Minute Squared 1 Femtometer/Second Squared in Meter/Minute Squared is equal to 3.6e-12 1 Femtometer/Second Squared = 3.6e-9 Millimeter/Minute Squared 1 Femtometer/Second Squared in Millimeter/Minute Squared is equal to 3.6e-9 1 Femtometer/Second Squared = 3.6e-10 Centimeter/Minute Squared 1 Femtometer/Second Squared in Centimeter/Minute Squared is equal to 3.6e-10 1 Femtometer/Second Squared = 3.6e-15 Kilometer/Minute Squared 1 Femtometer/Second Squared in Kilometer/Minute Squared is equal to 3.6e-15 1 Femtometer/Second Squared = 3.6e-15 Kilometer/Hour/Second 1 Femtometer/Second Squared in Kilometer/Hour/Second is equal to 3.6e-15 1 Femtometer/Second Squared = 8.503937007874e-9 Inch/Hour/Minute 1 Femtometer/Second Squared in Inch/Hour/Minute is equal to 8.503937007874e-9 1 Femtometer/Second Squared = 1.4173228346457e-10 Inch/Hour/Second 1 Femtometer/Second Squared in Inch/Hour/Second is equal to 1.4173228346457e-10 1 Femtometer/Second Squared = 2.3622047244094e-12 Inch/Minute/Second 1 Femtometer/Second Squared in Inch/Minute/Second is equal to 2.3622047244094e-12 1 Femtometer/Second Squared = 5.1023622047244e-7 Inch/Hour Squared 1 Femtometer/Second Squared in Inch/Hour Squared is equal to 5.1023622047244e-7 1 Femtometer/Second Squared = 1.4173228346457e-10 Inch/Minute Squared 1 Femtometer/Second Squared in Inch/Minute Squared is equal to 1.4173228346457e-10 1 Femtometer/Second Squared = 3.9370078740157e-14 Inch/Second Squared 1 Femtometer/Second Squared in Inch/Second Squared is equal to 3.9370078740157e-14 1 Femtometer/Second Squared = 7.0866141732283e-10 Feet/Hour/Minute 1 Femtometer/Second Squared in Feet/Hour/Minute is equal to 7.0866141732283e-10 1 Femtometer/Second Squared = 1.1811023622047e-11 Feet/Hour/Second 1 Femtometer/Second Squared in Feet/Hour/Second is equal to 1.1811023622047e-11 1 Femtometer/Second Squared = 1.9685039370079e-13 Feet/Minute/Second 1 Femtometer/Second Squared in Feet/Minute/Second is equal to 1.9685039370079e-13 1 Femtometer/Second Squared = 4.251968503937e-8 Feet/Hour Squared 1 Femtometer/Second Squared in Feet/Hour Squared is equal to 4.251968503937e-8 1 Femtometer/Second Squared = 1.1811023622047e-11 Feet/Minute Squared 1 Femtometer/Second Squared in Feet/Minute Squared is equal to 1.1811023622047e-11 1 Femtometer/Second Squared = 3.2808398950131e-15 Feet/Second Squared 1 Femtometer/Second Squared in Feet/Second Squared is equal to 3.2808398950131e-15 1 Femtometer/Second Squared = 6.9978402e-12 Knot/Hour 1 Femtometer/Second Squared in Knot/Hour is equal to 6.9978402e-12 1 Femtometer/Second Squared = 1.1663067e-13 Knot/Minute 1 Femtometer/Second Squared in Knot/Minute is equal to 1.1663067e-13 1 Femtometer/Second Squared = 1.9438445e-15 Knot/Second 1 Femtometer/Second Squared in Knot/Second is equal to 1.9438445e-15 1 Femtometer/Second Squared = 1.9438445e-18 Knot/Millisecond 1 Femtometer/Second Squared in Knot/Millisecond is equal to 1.9438445e-18 1 Femtometer/Second Squared = 1.3421617752326e-13 Mile/Hour/Minute 1 Femtometer/Second Squared in Mile/Hour/Minute is equal to 1.3421617752326e-13 1 Femtometer/Second Squared = 2.2369362920544e-15 Mile/Hour/Second 1 Femtometer/Second Squared in Mile/Hour/Second is equal to 2.2369362920544e-15 1 Femtometer/Second Squared = 8.0529706513958e-12 Mile/Hour Squared 1 Femtometer/Second Squared in Mile/Hour Squared is equal to 8.0529706513958e-12 1 Femtometer/Second Squared = 2.2369362920544e-15 Mile/Minute Squared 1 Femtometer/Second Squared in Mile/Minute Squared is equal to 2.2369362920544e-15 1 Femtometer/Second Squared = 6.2137119223733e-19 Mile/Second Squared 1 Femtometer/Second Squared in Mile/Second Squared is equal to 6.2137119223733e-19 1 Femtometer/Second Squared = 1.0936132983377e-15 Yard/Second Squared 1 Femtometer/Second Squared in Yard/Second Squared is equal to 1.0936132983377e-15 1 Femtometer/Second Squared = 1e-13 Gal 1 Femtometer/Second Squared in Gal is equal to 1e-13 1 Femtometer/Second Squared = 1e-13 Galileo 1 Femtometer/Second Squared in Galileo is equal to 1e-13 1 Femtometer/Second Squared = 1e-11 Centigal 1 Femtometer/Second Squared in Centigal is equal to 1e-11 1 Femtometer/Second Squared = 1e-12 Decigal 1 Femtometer/Second Squared in Decigal is equal to 1e-12 1 Femtometer/Second Squared = 1.0197162129779e-16 G-unit 1 Femtometer/Second Squared in G-unit is equal to 1.0197162129779e-16 1 Femtometer/Second Squared = 1.0197162129779e-16 Gn 1 Femtometer/Second Squared in Gn is equal to 1.0197162129779e-16 1 Femtometer/Second Squared = 1.0197162129779e-16 Gravity 1 Femtometer/Second Squared in Gravity is equal to 1.0197162129779e-16 1 Femtometer/Second Squared = 1e-10 Milligal 1 Femtometer/Second Squared in Milligal is equal to 1e-10 1 Femtometer/Second Squared = 1e-16 Kilogal 1 Femtometer/Second Squared in Kilogal is equal to 1e-16
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KEAM (Engineering) 2010 Questions (MCQ) on Electronics "You may never know what results come of your actions, but if you do nothing, there will be no results." –Mahatma Gandhi Questions on electronics will be generally interesting to most of you. Today we will discuss questions in this section which appeared in Kerala engineering entrance (KEAM - Engineering) 2010 question paper. Here are the questions with their solution: (1) A full wave rectifier with an a.c. input is shown: The output voltage across R[L] is represented as The rectified output voltage will be a direct voltage but there will be very large amount of ripples. The capacitor C acts as a filter to remove the ripples; but there will still be a small amount of ripples in the output. Therefore the correct option is (e). (2) In the given circuit the current through the battery is (b) 1 A (c) 1.5 A (d) 2 A (e) 2.5 A Since the diode D[1] is reverse biased, no current will flow through the D[1] branch. Diodes D[2] and D[3] are forward biased and hence the battery drives currents through the 20 Ω resistor and the series combination of the two 5 Ω resistors. The current driven through the 20 Ω resistor is 10 V/20 Ω = 0.5 A. The current driven through the 10 Ω resistor is 10 V/10 Ω = 1 A. Therefore, total current through the battery is 0.5 A + 1 A = 1.5 A (3) The collector supply voltage is 6 V and the voltage drop across a resistor of 600 Ω in the collector circuit is 0.6 V, in a transistor connected in common emitter mode. If the current gain is 20, the base current is (a) 0.25 mA (b) 0.05 mA (c) 0.12 mA (d) 0.02 mA (e) 0.07 mA We have I[C]R[C] = 0.6 V where I[C ]is the collector current and R[C] is the resistance in the collector circuit. Therefore, I[C]×600 Ω = 0.6 V from which I[C] = 0.6/600 A = 10^–3 A = 1 mA. Since the current gain β is given by β = I[C]/I[B] where I[B] is the base current, we have I[B] = I[C]/β = 1 mA/20 = 0.05 mA. (4) A pure semiconductor has equal electron and hole concentration of 10^16 m^–3. Doping by indium increases n[h] to 5×10^22 m^–3. Then the value of n[e] in the doped semiconductor is (a) 10^6 m^–3 (b) 10^22 m^–3 (c) 2×10^6 m^–3 (d) 10^19 m^–3 (e) 2×10^9 m^–3 According to the law of mass action we have n[i]^2 = n[e]n[h] where n[i] is the electron concentration as well as the hole concentration in the intrinsic (pure) semiconductor, n[e] is the electron concentration in the doped semiconductor and n [h] is the hole concentration in the doped semiconductor. Therefore n[e] = n[i]^2/n[h] = (10^16)^2/(5×10^22) = 2×10^9 m^–3
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About this project Biodose Tools is an open source project that aims to be a tool to perform all different tests and calculations needed by biological dosimetry laboratories. The app is developed using the R programming language and Shiny as a framework to offer an online, easy-to-use solution. Although the intention is to provide the application as a website, all R routines are available as an R package, which can be downloaded for improvement or personal use. We also aim to clarify and explain the tests used and to propose those considered most appropriate. Each laboratory in its routine work should choose the most suitable method, but the project aims to reach a consensus that will help us in case of mutual assistance or intercomparisons. The project is initially developed by RENEB association, but contributions are always welcome. Dicentrics: Dose-effect fitting Dicentrics: Dose estimation Curve fitting data options Variance-covariance matrix Translocations: Dose-effect fitting Genomic conversion factor Translocations: Dose estimation Curve fitting data options Variance-covariance matrix Genomic conversion factor Micronuclei: Dose-effect fitting Micronuclei: Dose estimation Curve fitting data options Variance-covariance matrix
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The algebraic structure of a trading stop-loss system I was once an undergraduate student in a joint Mathematics & Physics program. Some of the math courses, namely group theory and algebra, remained very abstract to me throughout my education. There is some group theory in the description of symmetries of physical systems; but being an experimentalist, I didn’t use more than 5% of what I learned in my undergrad during my PhD. However, in the course of my work now in finance, I had the pleasure of discovering that I was actually working with an algebraic structure. This post describes how that happened. The small trading firm for which I work is focusing a bit more on automated performance monitoring these days. With detailed trading performance data streaming in, it is now a good time to implement a stop-loss system. A stop-loss system is a system which receives trading performance data, and emits three categories of signal: • an all-clear signal, meaning that nothing in recent trading performance indicates a problem; • a warning signal, meaning that recent trading performance is degraded – but not yet concerning – and a human should take a look under the hood; • a halt signal, meaning that there is most probably something wrong, trading should be halted at once. Of course, we’re trading different products in different markets and even jurisdictions, and therefore the trading performance of every product is monitored independently. Moreover, our risk tolerance or expectations may be different for every product, and so a stop-loss system is really a framework in which to express multiple stop-loss rules, with different products being supervised by completely different stop-loss rules. Let us consider examples: assume that we’re trading a particular stock like AAPL^1. Sensible stop-loss rules might be: • If our current position has lost >10% in value over the last month, emit a warning; if the position has lost >25% over the last month, emit a halt signal. • If we’re expecting market volatility in the next hour to be high (for example, due to expected high-impact news), emit a halt signal. • If our forecast of the ticker price is way off – perhaps due to a problem in the forecasting model –, emit a halt signal. Here is what a rule framework might looks like^2: from enum import Enum, auto, unique from typing import Callable class Signal(Enum): AllClear = auto() Warn = auto() Halt = auto() class Context: Rule = Callable[[Context], Signal] # Example rule def rule(context: Context) -> Signal: A Rule is a function from some Context object to a Signal. We’re packing all information required to make decisions in a single data structure for reasons which will become obvious shortly. In this framework, we may express one of the stop loss rule examples as: def rule(context: Context) -> Signal: recent_loss = loss_percent( context.recent_performance(period="30d") ) if recent_loss > 0.25: return Signal.Halt elif recent_loss > 0.10 return Signal.Warn return Signal.AllClear For the remainder of this post, I don’t care anymore about the domain-specific content of a rule. My colleagues and I are expecting that, in practice, we will have pretty complex rules. In order to build complex rules from smaller, simpler rules, I wanted to be able to compose Rules together. This is straightforward because all rules have the same input and output types. Consider two rules, rule1 and rule2. If I want a new rule to halt if both rule1 and rule2 emit Signal.Halt, I could write it like this: def rule1(context: Context) -> Signal: def rule2(context: Context) -> Signal: def rule_lax(context: Context) -> Signal: sig1 = rule1(context) sig2 = rule2(context) if sig1 == sig2 == Signal.Halt: return Signal.Halt elif sig1 == sig2 == Signal.Warn: return Signal.Warn return Signal.AllClear That is an acceptable definition of rule composition. Since rule_lax will emit a Halt signal if both sub-rules emit a Halt signal, we’ll call this type of composition conjunction. In order to make it more ergonomic to write, let us wrap all rules in an object and re-use the & (overloaded and) operator: from dataclasses import dataclass from enum import Enum from operator import attrgetter class Signal(Enum): Signals can be composed using (&): >>> Signal.AllClear & Signal.AllClear < Signal.AllClear: 1 > >>> Signal.Warn & Signal.Halt < Signal.Warn: 2 > >>> Signal.Halt & Signal.Halt < Signal.Halt: 3 > AllClear = 1 Warn = 2 Halt = 3 def __and__(self, other: "Signal") -> "Signal": return min(self, other, key=attrgetter('value')) class rule(Callable): _inner: Callable[[Context], Signal] def __call__(self, context: Context) -> Signal: return self._inner.__call__(context=context) def __and__(self, other: "rule"): def newinner(context: Context) -> Signal: return rule1(context) & rule2(context) return self.__class__(newinner) and now we can re-write rule_lax like so: # The @rule decorator is required in order to lift rule1 from a regular function # to the `rule` object def rule1(context: Context) -> Signal: def rule2(context: Context) -> Signal: rule_lax = rule1 & rule2 Now, rule_lax is defined such that it’ll emit Signal.Halt if both rule1 and rule2 emit Signal.Halt. The same is true of warnings; if both rules emit a warning, then rule_lax will emit Signal.Warning. Here is a table which summarizes this composition: $A$ $B$ $A ~ \& ~ B$ $C$ $C$ $C$ $C$ $W$ $C$ $C$ $H$ $C$ $W$ $C$ $C$ $W$ $W$ $W$ $W$ $H$ $W$ $H$ $C$ $C$ $H$ $W$ $W$ $H$ $H$ $H$ where $C$ is Signal.AllClear, $W$ is Signal.Warning, and $H$ is Signal.Halt. Therefore, & is a binary function from Rules to Rule. This is not the only natural way to compose rules. What about this? def rule_strict(context: Context) -> Signal: sig1 = rule1(context) sig2 = rule2(context) if (sig1 == Signal.Halt) or (sig2 == Signal.Halt): return Signal.Halt elif (sig1 == Signal.Warning) or (sig2 == Signal.Warning): return Signal.Warning return Signal.AllClear In this case, rule_strict is more, uh, strict than rule_lax; it emits Signal.Halt if either rule1 or rule2 emits a stop signal. We’ll call this composition disjunction and re-use the | (overloaded or) operator to make it more ergonomic to write: class Signal(Enum): Signals can be composed using (&) and (|): >>> Signal.AllClear & Signal.AllClear < Signal.AllClear: 1 > >>> Signal.Warn & Signal.Halt < Signal.Warn: 2 > >>> Signal.Warn | Signal.Halt < Signal.Halt: 3 > AllClear = 1 Warn = 2 Halt = 3 def __and__(self, other: "Signal") -> "Signal": return min(self, other, key=attrgetter('value')) def __or__(self, other: "Signal") -> "Signal": return max(self, other, key=attrgetter('value')) class rule(Callable): _inner: Callable[[Context], Signal] def __call__(self, context: Context) -> Signal: return self._inner.__call__(context=context) def __and__(self, other: "rule"): def newinner(context: Context) -> Signal: return rule1(context) & rule2(context) return self.__class__(newinner) def __or__(self, other: "rule"): def newinner(context: Context) -> Signal: return rule1(context) | rule2(context) return self.__class__(newinner) With this implementation, we can express rule_lax and rule_strict as: # The @rule decorator is required in order to lift rule1 from a regular function # to the `rule` object def rule1(context: Context) -> Signal: def rule2(context: Context) -> Signal: rule_lax = rule1 & rule2 rule_strict = rule1 | rule2 We can update the table for the definition of & and |: $A$ $B$ $A ~ \& ~ B$ $A ~ | ~ B$ $C$ $C$ $C$ $C$ $C$ $W$ $C$ $W$ $C$ $H$ $C$ $H$ $W$ $C$ $C$ $W$ $W$ $W$ $W$ $W$ $W$ $H$ $W$ $H$ $H$ $C$ $C$ $H$ $H$ $W$ $W$ $H$ $H$ $H$ $H$ $H$ So for a given a given Context, which is fixed when the trading stop-loss system is running, we have: • A set of rule outcomes of type Signal; • A binary operation called conjunction (the & operator); □ & is associative; □ & is commutative; □ & has an identity, Signal.Halt; □ & does NOT have an inverse element. • A binary operation called disjunction (the | operator). □ | is associative; □ | is commutative; □ | has an identity, Signal.AllClear; □ | does NOT have an inverse element. That looks like a commutative semiring to me! Just a few more things to check: • | distributes from both sides over &: □ $a ~|~ (b ~\&~ c)=(a ~|~ b) ~\&~ (a ~\&~ c)$ for all $a$, $b$, and $c$; □ $(a ~ \& ~ b) ~|~ c = (a ~|~ c) ~\&~ (b ~\&~ c)$ for all $a$, $b$, and $c$. • The identity element of & (called $0$, in this case Signal.Halt) annihilates the | operation, i.e. $0 ~ | ~ a = 0$ for all $a$. Don’t take my word for it, we can check exhaustively: from itertools import product zero = Signal.Halt one = Signal.AllClear # Assert & is associative assert all( (a & b) & c == a & (b & c) for (a, b, c) in product(Signal, repeat=3) ) # Assert & is commutative assert all( a & b == b & a for (a, b) in product(Signal, repeat=2) ) # Assert & has an identity assert all( a & zero == a for a in Signal ) # Assert | is associative assert all( (a | b) | c == a | (b | c) for (a, b, c) in product(Signal, repeat=3) ) # Assert | has an identity assert all( a | one == a for a in Signal ) # Assert | distributes over & on both sides assert all( a | (b & c) == (a | b) & (a | c) for (a, b, c) in product(Signal, repeat=3) ) assert all( (a & b) | c == (a | c) & (b | c) for (a, b, c) in product(Signal, repeat=3) ) # Assert identity of & annihilates with respect to | assert all( (zero | a) == zero for a in Signal) and there we have it! This design of a trading stop-loss system is an example of commutative semirings. This fact does absolutely nothing in the practical sense; I’m just happy to have spotted this structure more than 10 years after seeing it in undergrad. 1. I’m actually not involved in trading securities at all, but I think intuition about stock markets is more common↩︎ 2. I’ll be using Python in this post because it was a requirement of the implementation, but know that I’m doing this under protest.↩︎
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testing ibpdv and ibpu 05-12-2015, 01:57 PM Post: #1 salvomic Posts: 1,396 Senior Member Joined: Jan 2015 testing ibpdv and ibpu hi all, I'm testing the "integration by parts" in the Prime. Inputting ibpu(x*sin(x),x) I get [-x*cos(x) cos(x)], ok with ibpdv(x*sin(x),sin(x) I get [x*sin(x)^2/cos(x) (-x*sin(x)^3-cos(x)*sin(x)^2-x*cos(x)^2)/cos(x)^2] I'm a bit confused: shouldn't I get the same result as above? Am I using wrong syntax? :-) Then another question: there is also a "integration by substitution" method in the Prime? thank you ∫aL√0mic (IT9CLU) :: HP Prime 50g 41CX 71b 42s 39s 35s 12C 15C - DM42, DM41X - WP34s Prime Soft. Lib 05-12-2015, 02:44 PM Post: #2 parisse Posts: 1,337 Senior Member Joined: Dec 2013 RE: testing ibpdv and ibpu The second argument of ibpdv should be the antiderivative of one term in the product. 05-12-2015, 02:45 PM (This post was last modified: 05-12-2015 02:46 PM by salvomic.) Post: #3 salvomic Posts: 1,396 Senior Member Joined: Jan 2015 RE: testing ibpdv and ibpu (05-12-2015 02:44 PM)parisse Wrote: The second argument of ibpdv should be the antiderivative of one term in the product. thank you! in fact, doing so it works well! ibpdv(x*sin(x),-cos(x) -> [-x*cos(x) cos(x)] ∫aL√0mic (IT9CLU) :: HP Prime 50g 41CX 71b 42s 39s 35s 12C 15C - DM42, DM41X - WP34s Prime Soft. Lib 05-12-2015, 07:45 PM Post: #4 Aries Posts: 159 Member Joined: Oct 2014 RE: testing ibpdv and ibpu Hey Salvo, as to the integration by substitution, take a look here: 05-12-2015, 07:59 PM (This post was last modified: 05-12-2015 08:05 PM by salvomic.) Post: #5 salvomic Posts: 1,396 Senior Member Joined: Jan 2015 RE: testing ibpdv and ibpu (05-12-2015 07:45 PM)Aries Wrote: Hey Salvo, as to the integration by substitution, take a look here: thank you, Aries! I'm reading it and trying \[ \int_0^1{\sqrt{1-x^{2}}}dx \] assume(u>0); a:=subst('integrate(sqrt(1-x^2)),0,1,x)',x=sin(u)); a; and I get [u, π/4 π/4] Manually I have solution = (½)arcsin1 = π/4 (only?) Is it right? I thought to have this series: ∫cos(u)^2du, then (u/2)+(sin(2u))/4+c and finally (½)(arcsin(x)+x*sqrt(1-x^2)+c (changing also the integral bounds) ... However it's interesting this approach. ∫aL√0mic (IT9CLU) :: HP Prime 50g 41CX 71b 42s 39s 35s 12C 15C - DM42, DM41X - WP34s Prime Soft. Lib 05-12-2015, 10:18 PM (This post was last modified: 05-12-2015 10:29 PM by Arno K.) Post: #6 Arno K Posts: 472 Senior Member Joined: Mar 2015 RE: testing ibpdv and ibpu Hello Salvo, as I am fairly new to the Prime I spend a lot of time trying it, so I entered your Integration example and got the list :[u, int(sqrt(1-sin(u)^2)*cos(u),u,0,PI/2),PI/4]. This is: [substitution var, new integral, result], your solution is only shown after a final hit on simplify. The list is produced because you are entering three commands in one line. So :simplify (subst('int(sqrt(1-x^2),x,0,1)',x=sin(u))) directly produces nothing but PI/4. Now I tried int(sqrt(1-x^2),x) to see what the Prime would deliver: the answer is too long to be typed here,containing i and LN aside the first part, I thought of getting:0.5*(x*sqrt(1-x^2)+arcsin So I entered: assume(x,float) and gave it a try, no change, I added: additionally(x>=0), still no change. By the way, Complex is switched off. Perhaps you or someone else can give it a try. edit: I now took the real Prime and entered the integral which clearly produced the same result, xcas on my Computer, xcas pad on my mobile and my formerly preferred cas (MathStudio for Android) delivered the simplified answer, which you provided and I had expected. I am a bit disappointed... 05-13-2015, 05:14 AM Post: #7 parisse Posts: 1,337 Senior Member Joined: Dec 2013 RE: testing ibpdv and ibpu (05-12-2015 07:59 PM)salvomic Wrote: I'm reading it and trying \[ \int_0^1{\sqrt{1-x^{2}}}dx \] assume(u>0); a:=subst('integrate(sqrt(1-x^2)),0,1,x)',x=sin(u)); a; It should be assume(u>0); a:=subst('integrate(sqrt(1-x^2),x,0,1)',x=sin(u)); a 05-13-2015, 07:18 AM Post: #8 salvomic Posts: 1,396 Senior Member Joined: Jan 2015 RE: testing ibpdv and ibpu (05-12-2015 10:18 PM)Arno K Wrote: ... So :simplify(subst('int(sqrt(1-x^2),x,0,1)',x=sin(u))) directly produces nothing but PI/4. Now I tried int(sqrt(1-x^2),x) to see what the Prime would deliver: the answer is too long to be typed here,containing i and LN aside the first part, I am a bit disappointed... (05-13-2015 05:14 AM)parisse Wrote: It should be assume(u>0); a:=subst('integrate(sqrt(1-x^2),x,0,1)',x=sin(u)); a thank you Parisse! iw would be interesting also to have a method to get "step by step" solution, here and for the "integration by parts" functions, but anyway still already this is ok... please, control your CAS Settings: I've in CAS: Exact, "Use √" and Principal checked, Complex not checked (with firmware 6975) with the integral you should get (½)*ASIN(x)+(½)*x*√(-x^2-1) from 0 to 1, that's π/4 (in Home you gets 0.785398...) try, and let me know... ∫aL√0mic (IT9CLU) :: HP Prime 50g 41CX 71b 42s 39s 35s 12C 15C - DM42, DM41X - WP34s Prime Soft. Lib 05-13-2015, 07:34 AM Post: #9 Arno K Posts: 472 Senior Member Joined: Mar 2015 RE: testing ibpdv and ibpu Yes, the bounded integral is no problem, the one I tried later , int(sqrt(1-x^2),x) without boundaries, does not work. 05-13-2015, 07:57 AM Post: #10 salvomic Posts: 1,396 Senior Member Joined: Jan 2015 RE: testing ibpdv and ibpu (05-13-2015 07:34 AM)Arno K Wrote: Yes, the bounded integral is no problem, the one I tried later , int(sqrt(1-x^2),x) without boundaries, does not work. but have you already controlled the settings? perhaps it's that... I'm using textbook. Otherwise, I don't know I've the correct result also in emulator (firmware 6975) ∫aL√0mic (IT9CLU) :: HP Prime 50g 41CX 71b 42s 39s 35s 12C 15C - DM42, DM41X - WP34s Prime Soft. Lib 05-13-2015, 08:16 AM Post: #11 Arno K Posts: 472 Senior Member Joined: Mar 2015 RE: testing ibpdv and ibpu Hi Salvo, that are the same cas-Settings as mine, , Software Version is 6975 on the emulator as well as on my calculator, I am not able to get the symbolic integral, that is PI/4 is no problem 05-13-2015, 08:33 AM Post: #12 salvomic Posts: 1,396 Senior Member Joined: Jan 2015 RE: testing ibpdv and ibpu (05-13-2015 08:16 AM)Arno K Wrote: Hi Salvo, that are the same cas-Settings as mine, , Software Version is 6975 on the emulator as well as on my calculator, I am not able to get the symbolic integral, that is PI/4 is no problem I think however it's only a settings' problem. Let's see... CAS settings: simplify "minimum"; Home settings Entry "textbook", complex "a+b*i", allow complex output... ok (but those shouldn't change much)... Use helper on key C: insert ∫ and so on (without boundaries), I can see the graphic (textbook) integral, and after Enter I get the correct result. I'm in CAS mode, the "x variable" is x and it isn't "valued" (check that it doesn't have value or use use purge(x)). It *should* work by you ∫aL√0mic (IT9CLU) :: HP Prime 50g 41CX 71b 42s 39s 35s 12C 15C - DM42, DM41X - WP34s Prime Soft. Lib 05-13-2015, 08:52 AM (This post was last modified: 05-13-2015 09:17 AM by Arno K.) Post: #13 Arno K Posts: 472 Senior Member Joined: Mar 2015 RE: testing ibpdv and ibpu Yes, you are right, I reset the Emulator, entered the integral and it worked fine. Now I reinstalled my backup, reentered the integral and got the same weird result than before. Your CAS Settings don't work either, the helper from key C doesn´t change anything. So it is up to me, to look for a solution, the Problem seems to be caused by one of the installed programs which I keep on the same level on both the emu and the calc, I now start thinking of having the emu empty for testing and then using these things on the calc edit: after playing around with various settings I again reset the Emulator, then looked at the settings there, reinstalled the backup, set everything like in the fresh emu and it worked. The whole problem is caused by: use i, which I had checked in the cas-settings. I hadn't awaited that, this flag was, in my opinion, for the calc to be used for factorization of polynomials (I don't see the point where here a polynomial is factored but perhaps can M. Parisse help and make things clearer) and perhaps for output purposes of complex numbers , i.e. (3,5) versus 3+5*i, but it seems to influence other calculations as well. It would be nice if things like that, which are rather annoying, were clearly mentioned in the manual. 05-13-2015, 08:57 AM Post: #14 salvomic Posts: 1,396 Senior Member Joined: Jan 2015 RE: testing ibpdv and ibpu (05-13-2015 08:52 AM)Arno K Wrote: Yes, you are right, I reset the Emulator, entered the integral and it worked fine. Now I reinstalled my backup, reentered the integral and got the same weird result than before. Your CAS Settings don't work either, the helper from key C doesn´t change anything. So it is up to me, to look for a solution, the Problem seems to be caused by one of the installed programs which I keep on the same level on both the emu and the calc, I now start thinking of having the emu empty for testing and then using these things on the calc Have already compared the Settings (both CAS and Home) in Emulator (after reset) and the Prime? Maybe there is one different... I don't know if a custom program could change the result of the integral, I think this is a bit strange... Try again and then tell us, I'm curious about this weird result. ∫aL√0mic (IT9CLU) :: HP Prime 50g 41CX 71b 42s 39s 35s 12C 15C - DM42, DM41X - WP34s Prime Soft. Lib 05-13-2015, 09:32 AM Post: #15 Arno K Posts: 472 Senior Member Joined: Mar 2015 RE: testing ibpdv and ibpu (05-13-2015 08:57 AM)salvomic Wrote: (05-13-2015 08:52 AM)Arno K Wrote: Yes, you are right, I reset the Emulator, entered the integral and it worked fine. Now I reinstalled my backup, reentered the integral and got the same weird result than before. Your CAS Settings don't work either, the helper from key C doesn´t change anything. So it is up to me, to look for a solution, the Problem seems to be caused by one of the installed programs which I keep on the same level on both the emu and the calc, I now start thinking of having the emu empty for testing and then using these things on the calc Have already compared the Settings (both CAS and Home) in Emulator (after reset) and the Prime? Maybe there is one different... I don't know if a custom program could change the result of the integral, I think this is a bit strange... Try again and then tell us, I'm curious about this weird result. see my edit above 05-13-2015, 09:43 AM Post: #16 salvomic Posts: 1,396 Senior Member Joined: Jan 2015 RE: testing ibpdv and ibpu (05-13-2015 09:32 AM)Arno K Wrote: see my edit above I read! yes, I tried here: "use i" produce your same output, using "i" evaluating the integral, unchecking it, it works fine. ∫aL√0mic (IT9CLU) :: HP Prime 50g 41CX 71b 42s 39s 35s 12C 15C - DM42, DM41X - WP34s Prime Soft. Lib 05-13-2015, 02:43 PM Post: #17 parisse Posts: 1,337 Senior Member Joined: Dec 2013 RE: testing ibpdv and ibpu If you check use i (Xcas complex mode), you will get complex logarithms instead of inverse trigonometric functions. 05-13-2015, 05:22 PM Post: #18 Arno K Posts: 472 Senior Member Joined: Mar 2015 RE: testing ibpdv and ibpu Hello parisse, yes, I noticed this behaviour as you can see above, but the help on this flag which can be seen in CAS-Settings claims something like: Use i to factor polynomials and not that it influences production on various kinds of results, so I thought it would not matter if switched on or off. Now I know better. Thanks Arno User(s) browsing this thread: 1 Guest(s)
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Generate Untargeted and Targeted Adversarial Examples for Image Classification- MATLAB & Simulink- MathWorks 日本 (2024) This example shows how to use the fast gradient sign method (FGSM) and the basic iterative method (BIM) to generate adversarial examples for a pretrained neural network. Neural networks can be susceptible to a phenomenon known as adversarial examples [1], where very small changes to an input can cause the input to be misclassified. These changes are often imperceptible to humans. In this example, you create two types of adversarial examples: • Untargeted — Modify an image so that it is misclassified as any incorrect class. • Targeted — Modify an image so that it is misclassified as a specific class. Load Network and Image Load a network that has been trained on the ImageNet [2] data set and convert it to a dlnetwork. net = squeezenet;lgraph = layerGraph(net);lgraph = removeLayers(lgraph,lgraph.Layers(end).Name);dlnet = dlnetwork(lgraph); Extract the class labels. classes = categories(net.Layers(end).Classes); Load an image to use to generate an adversarial example. The image is a picture of a golden retriever. img = imread('sherlock.jpg');T = "golden retriever"; Resize the image to match the input size of the network. inputSize = dlnet.Layers(1).InputSize;img = imresize(img,inputSize(1:2));figureimshow(img)title("Ground Truth: " + T) Prepare the image by converting it to a dlarray. X = dlarray(single(img),"SSCB"); Prepare the label by one-hot encoding it. T = onehotencode(T,1,'ClassNames',classes);T = dlarray(single(T),"CB"); Untargeted Fast Gradient Sign Method Create an adversarial example using the untargeted FGSM [3]. This method calculates the gradient ${abla }_{X}L\left(X,T\right)$ of the loss function $L$, with respect to the image $X$ you want to find an adversarial example for, and the class label $T$. This gradient describes the direction to "push" the image in to increase the chance it is misclassified. You can then add or subtract a small error from each pixel to increase the likelihood the image is misclassified. The adversarial example is calculated as follows: ${\mathit{X}}_{\mathrm{adv}}=\mathit{X}+ϵ.\mathrm{sign}\left({abla }_{\mathit{X}}\mathit{L}\left(\mathit{X},\mathit{T}\right)\right)$. Parameter $ϵ$ controls the size of the push. A larger $ϵ$ value increases the chance of generating a misclassified image, but makes the change in the image more visible. This method is untargeted, as the aim is to get the image misclassified, regardless of which class. Calculate the gradient of the image with respect to the golden retriever class. gradient = dlfeval(@untargetedGradients,dlnet,X,T); Set epsilon to 1 and generate the adversarial example. epsilon = 1;XAdv = X + epsilon*sign(gradient); Predict the class of the original image and the adversarial image. YPred = predict(dlnet,X);YPred = onehotdecode(squeeze(YPred),classes,1) YPred = categorical golden retriever YPredAdv = predict(dlnet,XAdv);YPredAdv = onehotdecode(squeeze(YPredAdv),classes,1) YPredAdv = categorical Labrador retriever Display the original image, the perturbation added to the image, and the adversarial image. If the epsilon value is large enough, the adversarial image has a different class label from the original The network correctly classifies the unaltered image as a golden retriever. However, because of perturbation, the network misclassifies the adversarial image as a labrador retriever. Once added to the image, the perturbation is imperceptible, demonstrating how adversarial examples can exploit robustness issues within a network. Targeted Adversarial Examples A simple improvement to FGSM is to perform multiple iterations. This approach is known as the basic iterative method (BIM) [4] or projected gradient descent [5]. For the BIM, the size of the perturbation is controlled by parameter $\alpha$ representing the step size in each iteration. This is as the BIM usually takes many, smaller, FGSM steps in the direction of the gradient. After each iteration, clip the perturbation to ensure the magnitude does not exceed $ϵ$. This method can yield adversarial examples with less distortion than FGSM. When you use untargeted FGSM, the predicted label of the adversarial example can be very similar to the label of the original image. For example, a dog might be misclassified as a different kind of dog. However, you can easily modify these methods to misclassify an image as a specific class. Instead of maximizing the cross-entropy loss, you can minimize the mean squared error between the output of the network and the desired target output. Generate a targeted adversarial example using the BIM and the great white shark target class. targetClass = "great white shark";targetClass = onehotencode(targetClass,1,'ClassNames',classes); Increase the epsilon value to 5, set the step size alpha to 0.2, and perform 25 iterations. Note that you may have to adjust these settings for other networks. epsilon = 5;alpha = 0.2;numIterations = 25; Keep track of the perturbation and clip any values that exceed epsilon. delta = zeros(size(X),'like',X);for i = 1:numIterations gradient = dlfeval(@targetedGradients,dlnet,X+delta,targetClass); delta = delta - alpha*sign(gradient); delta(delta > epsilon) = epsilon; delta(delta < -epsilon) = -epsilon;endXAdvTarget = X + delta; Predict the class of the targeted adversarial example. YPredAdvTarget = predict(dlnet,XAdvTarget);YPredAdvTarget = onehotdecode(squeeze(YPredAdvTarget),classes,1) YPredAdvTarget = categorical great white shark Display the original image, the perturbation added to the image, and the targeted adversarial image. Because of imperceptible perturbation, the network classifies the adversarial image as a great white shark. To make the network more robust against adversarial examples, you can use adversarial training. For an example showing how to train a network robust to adversarial examples, see Train Image Classification Network Robust to Adversarial Examples. Supporting Functions Untargeted Input Gradient Function Calculate the gradient used to create an untargeted adversarial example. This gradient is the gradient of the cross-entropy loss. function gradient = untargetedGradients(dlnet,X,target)Y = predict(dlnet,X);Y = stripdims(squeeze(Y));loss = crossentropy(Y,target,'DataFormat','CB');gradient = dlgradient(loss,X);end Targeted Input Gradient Function Calculate the gradient used to create a targeted adversarial example. This gradient is the gradient of the mean squared error. function gradient = targetedGradients(dlnet,X,target)Y = predict(dlnet,X);Y = stripdims(squeeze(Y));loss = mse(Y,target,'DataFormat','CB');gradient = dlgradient(loss,X);end Show Adversarial Image Show an image, the corresponding adversarial image, and the difference between the two (perturbation). function showAdversarialImage(image,label,imageAdv,labelAdv,epsilon)figuresubplot(1,3,1)imgTrue = uint8(extractdata(image));imshow(imgTrue)title("Original Image" + newline + "Class: " + string(label))subplot(1,3,2)perturbation = uint8(extractdata(imageAdv-image+127.5));imshow(perturbation)title("Perturbation")subplot(1,3,3)advImg = uint8(extractdata(imageAdv));imshow(advImg)title("Adversarial Image (Epsilon = " + string(epsilon) + ")" + newline + ... "Class: " + string(labelAdv))end [1] Goodfellow, Ian J., Jonathon Shlens, and Christian Szegedy. “Explaining and Harnessing Adversarial Examples.” Preprint, submitted March 20, 2015. https://arxiv.org/abs/1412.6572. [2] ImageNet. http://www.image-net.org. [3] Szegedy, Christian, Wojciech Zaremba, Ilya Sutskever, Joan Bruna, Dumitru Erhan, Ian Goodfellow, and Rob Fergus. “Intriguing Properties of Neural Networks.” Preprint, submitted February 19, 2014. [4] Kurakin, Alexey, Ian Goodfellow, and Samy Bengio. “Adversarial Examples in the Physical World.” Preprint, submitted February 10, 2017. https://arxiv.org/abs/1607.02533. [5] Madry, Aleksander, Aleksandar Makelov, Ludwig Schmidt, Dimitris Tsipras, and Adrian Vladu. “Towards Deep Learning Models Resistant to Adversarial Attacks.” Preprint, submitted September 4, 2019. See Also dlnetwork | onehotdecode | onehotencode | predict | dlfeval | dlgradient | estimateNetworkOutputBounds | verifyNetworkRobustness Related Topics • Verification of Neural Networks • Train Image Classification Network Robust to Adversarial Examples • Generate Adversarial Examples for Semantic Segmentation • Grad-CAM Reveals the Why Behind Deep Learning Decisions • Understand Network Predictions Using LIME
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Pasteur and parachutes: when statistical process control is better than a randomized controlled trial Heroes and martyrs of quality and safety Pasteur and parachutes: when statistical process control is better than a randomized controlled trial Statistics from Altmetric.com Request Permissions If you wish to reuse any or all of this article please use the link below which will take you to the Copyright Clearance Center’s RightsLink service. You will be able to get a quick price and instant permission to reuse the content in many different ways. On 6 July 1885 Joseph Meister, a 9 year old boy who had been severely bitten 2 days before by a rabid dog, was treated in Paris with the rabies vaccine developed in Louis Pasteur’s (1822–95) laboratory after years of brilliant scientific research and experimentation on animals.^1,^2 Before this, no one who had developed the symptoms of rabies had survived. Meister was at the highest risk of developing symptoms but, after 10 days of vaccinations, he was fine and lived for many years. The results of this and a second case were so dramatic compared with previous experience that, by October, speakers at the French Academy of Sciences stated that it was “necessary to organize this treatment for everyone” and “This is a memorable day in the history of medicine”. Philanthropic contributions poured in and by 1888 the Pasteur Institute was founded. By then about 1200 patients had been vaccinated with a mortality rate of 1%.^3 Transmitted to humans by animal bites, rabies has always been a rare event. A biting animal may not have rabies. Verification that an animal is rabid is not always possible if the animal cannot be caught and watched. A rabid animal that bites does not always transmit the rabies virus. If the person is infected, the incubation period is about 20–60 days before symptoms develop leading to a painful and certain death.^4 Today there are occasional reports of survival but these are so rare as to make headlines.^5 The rarity of rabies and the long incubation period led to Pasteur’s novel approach of not vaccinating everyone preventively. The incubation period allowed time for his 10 days of vaccination to build immunity. Would you—having been infected by a rabid dog—be willing to participate in a randomized controlled trial (RCT) when being in the control group had a certainty of a “most awful death”? If volunteers could be found, the trial would have to be small and would therefore have low statistical power. If one wished to show that vaccination was effective for men and women—young, middle aged, and old—the sample size would have to be 264, condemning half those people to a certain death. In this example the application of statistical process control (SPC) makes more sense. SPC avoids the ethical issues, saves lives, builds on prior experience to control confounding variables, gives an answer more rapidly, and has much more statistical power. In order to demonstrate the application of statistical process control (SPC) by the use of control charts, we have simulated pre and post 1885 rabies mortality using data from the literature. This simulation assumes a mean (SD) survival of 20 (10) days before 1885. Patients were grouped into blocks of five so that each point on the control chart represents five patients recorded sequentially over time. This would be equivalent to an average of five patients bitten in a day. The survival in days is plotted for 500 groups of five patients sequentially. Survival in days after receiving Pasteur’s treatment is simulated for a similar number of patients based on a mean survival of 45 years combined with a mortality rate of 1% from the treatment (figs 1 and 2). The results are so dramatic that we have presented them in two formats. In fig 1 the vertical axis is measured in number of days of survival. In this presentation the variation before 1885 is so small as to be unobservable. Even the 3-sigma upper and lower control limits are too close together to be seen. The shorter survival points in the post 1885 treatment experience reflect one death in that group of five patients. In order to show more clearly the pre 1885 variation, fig 2 presents exactly the same data but with the vertical axis transformed into a logarithmic scale. The pre 1885 data show a very stable process without special cause variation. Every one with rabies symptoms soon died. The serious scholar may disagree with our simple approach and assumptions about survival, but we think that the before and after differences were so great that this simulation is plausible. The statistical power of this evidence is overwhelming. Based on these pre 1885 control limits, the probability of living 4000 days is above the 894-sigma level yet the results were even more dramatic than that. When little Joseph Meister had lived to October (or 90 days), his survival was at the 20-sigma level (p=6.4×10^−13) and the French academicians were right in declaring Pasteur’s treatment a great victory. Thus, SPC methods can demonstrate a dramatic difference as a result of the outcome from one patient and the new treatment can be started immediately, rather than waiting for the results of a prospective controlled trial. Gordon Smith and Jill Pell wrote a fine satirical article in 2003 pointing out that the use of parachutes has never been subject to a randomized controlled trial.^6 A careful literature review found rare examples of people falling from great heights and living. Olympic ski jumpers survive their falls. The authors report that parachutes can fail to open resulting in death. We all learn at an early age the unvarying power of gravity and do not need to be convinced that falling from a great height is likely to be fatal. SPC evaluation using past experience which we applied to rabies vaccination may be relevant here. Even more dramatic control charts could be simulated, particularly if ski jumpers are excluded. There are other examples. The progress made in the treatment of acute lymphoblastic leukemia (ALL) of childhood is one of the true success stories of modern medicine.^7 Incremental advances over 50 years mean that ALL has gone from a uniformly fatal disease to one with an overall cure rate of more than 75%. Another example was the dramatic introduction of ether as an anaesthetic during surgery. In this case the outcome measure would be pain rather than mortality. For the four dramatic improvements described here, we propose that information from prior experience using SPC is to be preferred to RCTs for six reasons (table 1). In these circumstances SPC has greater statistical power to exclude chance as an explanation. The RCT is designed to control for unknown confounding variables. Perhaps the treatment only works for young boys and not for older women. In the case of symptomatic rabies before 1885, men and women (young and old) all died without variation. If there were any unknown confounding variables they would appear as special cause variation in a long series of prior observations. SPC can give a very rapid answer in these circumstances. There needs to be a plausible process (treatment) change associated with the astonishing outcome that is replicable for the next patients. Without this scientifically based replicable treatment, Joseph Meister’s survival could be declared a miracle and his cure associated with the intervention of a saint—the statistical analysis of miracles. One test of an ethical randomized trial is whether you yourself and others would volunteer for random assignment to control or experimental groups. When the expected differences are small, when unknown confounding variables are likely to overwhelm the treatment effect, where the casual model is weak, where prior information is thin, when there is a single intervention and single end point, and results are not urgent—then RCTs are more useful. Appendix: Technical Note This section describes the simulation we performed to demonstrate the applicability of control charts to detect extraordinary causes in the rabies case study. We generated two samples to represent the survival of two populations of individuals after onset of rabies symptoms. They are before and after July 1885, the date in which Pasteur incidentally initiated the test of the rabies vaccine in humans. We modeled each population’s survival assuming that the survival time presents a Weibull distribution, a common probability distribution used to model time-to-an-event variables.^8 This distribution allows for a dependence of the hazard on time. In this case, this would represent a potential change in the risk of death with time since presentation of rabies symptoms. Also, in practical terms, this distribution has more flexibility for modeling than a simpler exponential distribution, also commonly used in survival analysis. Following is the survival function of a random variable exhibiting a Weibull distribution with parameters λ and γ representing the scale and shape parameters of this distribution: We selected for the first sample, before July 1885, the parameters λ and γ such that the mean and standard deviation of the survival time were close to 20 and 10 days, respectively. For the second sample, after July 1885, we chose the parameters λ and γ in order to represent survival time with mean and standard deviation of 45 and 5 years, respectively. This was the average survival time in France for 1885.^9 For this second sample we incorporated individuals in whom the vaccine did not work with a frequency of 1/100 individuals. We assumed for these particular patients a survival time identical to those for people before July 1885. For both samples we assumed a period of 15 months each. This was the interval reported in the study by Moulin^10 in which 2500 people had received the rabies vaccine. We assumed this rate of approximately five bitten people/day for both the before and after July 1885 samples. Table 2 shows the summary of the parameter values we used to generate both samples. Control charts are a key component of the area of statistical process control (SPC) initially developed and more disseminated in quality control of industrial processes. Several types of control charts are available depending on the nature of the measurement to control and the particular statistical control scheme. One of the purposes of control charts is to identify instances in which a particular process monitored through a particular measurement goes “out of control” and the potential causes for such situation. Another purpose is to improve the process itself. One of the most popular ones is the mean control chart developed by Shewart.^11 As its name indicates, the mean control chart displays the time evolution of the mean of a continuous variable of interest (fig 3). The chart has a distinctive pattern marked by three reference lines. One in the center, designated the center line (CL), set at the average value of the characteristic to monitor while “in control”. This value can be obtained from previous experience or assigned as a milestone. The other two reference lines assigned upper and lower control limits (UCL and LCL respectively) corresponding to the boundaries beyond which the process will be signaled as out of control. For mean control charts these limits are commonly set at three standard error of the process mean of the process in control or what is known as the 3-sigma limit in the SPC jargon. We performed all the computations for the simulation in SAS for Windows Version 9.1.
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Source code for dclab.lme4.wrapr """R lme4 wrapper""" import logging import numbers import pathlib import tempfile import importlib_resources import numpy as np from .. import definitions as dfn from ..rtdc_dataset.core import RTDCBase from . import rsetup logger = logging.getLogger(__name__) class Rlme4(object): def __init__(self, model="lmer", feature="deform"): """Perform an R-lme4 analysis with RT-DC data model: str One of: - "lmer": linear mixed model using lme4's ``lmer`` - "glmer+loglink": generalized linear mixed model using lme4's ``glmer`` with an additional a log-link function via the ``family=Gamma(link='log'))`` keyword. feature: str Dclab feature for which to compute the model #: modeling method to use (e.g. "lmer") self.model = None #: dclab feature for which to perform the analysis self.feature = None #: list of [RTDCBase, column, repetition, chip_region] self.data = [] self.set_options(model=model, feature=feature) # Make sure that lme4 is available if not rsetup.has_lme4(): logger.info("Installing lme4, this may take a while!") def add_dataset(self, ds, group, repetition): """Add a dataset to the analysis list ds: RTDCBase group: str The group the measurement belongs to ("control" or repetition: int Repetition of the measurement - For each repetition, there must be a "treatment" (``1``) and a "control" (``0``) group. - If you would like to perform a differential feature analysis, then you need to pass at least a reservoir and a channel dataset (with same parameters for `group` and `repetition`). assert group in ["treatment", "control"] assert isinstance(ds, RTDCBase) assert isinstance(repetition, numbers.Integral) region = ds.config["setup"]["chip region"] # make sure there are no doublets for ii, dd in enumerate(self.data): if dd[1] == group and dd[2] == repetition and dd[3] == region: raise ValueError("A dataset with group '{}', ".format(group) + "repetition '{}', and ".format(repetition) + "'{}' region has already ".format(region) + "been added (index {})!".format(ii)) self.data.append([ds, group, repetition, region]) def check_data(self): """Perform sanity checks on ``self.data``""" # Check that we have enough data if len(self.data) < 3: msg = "Linear mixed effects models require repeated " + "measurements. Please add more repetitions." raise ValueError(msg) def fit(self, model=None, feature=None): """Perform (generalized) linear mixed-effects model fit The response variable is modeled using two linear mixed effect - model: "feature ~ group + (1 + group | repetition)" (random intercept + random slope model) - the null model: "feature ~ (1 + group | repetition)" (without the fixed effect introduced by the "treatment" group). Both models are compared in R using "anova" (from the R-package "stats" :cite:`Everitt1992`) which performs a likelihood ratio test to obtain the p-Value for the significance of the fixed effect (treatment). If the input datasets contain data from the "reservoir" region, then the analysis is performed for the differential model: str (optional) One of: - "lmer": linear mixed model using lme4's ``lmer`` - "glmer+loglink": generalized linear mixed model using lme4's ``glmer`` with an additional log-link function via ``family=Gamma(link='log'))`` :cite:`lme4` feature: str (optional) dclab feature for which to compute the model results: dict Dictionary with the results of the fitting process: - "anova p-value": Anova likelihood ratio test (significance) - "feature": name of the feature used for the analysis - "fixed effects intercept": Mean of ``self.feature`` for all controls; In the case of the "glmer+loglink" model, the intercept is already back transformed from log space. - "fixed effects treatment": The fixed effect size between the mean of the controls and the mean of the treatments relative to "fixed effects intercept"; In the case of the "glmer+loglink" model, the fixed effect is already back transformed from log - "fixed effects repetitions": The effects (intercept and treatment) for each repetition. The first axis defines intercept/treatment; the second axis enumerates the repetitions; thus the shape is (2, number of repetitions) and ``np.mean(results["fixed effects repetitions"], axis=1)`` is equivalent to the tuple (``results["fixed effects intercept"]``, ``results["fixed effects treatment"]``) for the "lmer" model. This does not hold for the "glmer+loglink" model, because of the non-linear inverse transform back from log space. - "is differential": Boolean indicating whether or not the analysis was performed for the differential (bootstrapped and subtracted reservoir from channel data) feature - "model": model name used for the analysis ``self.model`` - "model converged": boolean indicating whether the model - "r model summary": Summary of the model - "r model coefficients": Model coefficient table - "r script": the R script used - "r output": full output of the R script self.set_options(model=model, feature=feature) # Assemble dataset if self.is_differential(): # bootstrap and compute differential features using reservoir features, groups, repetitions = self.get_differential_dataset() # regular feature analysis features = [] groups = [] repetitions = [] for dd in self.data: features.append(self.get_feature_data(dd[1], dd[2])) # concatenate and populate arrays for R features_c = np.concatenate(features) groups_c = np.zeros(len(features_c), dtype=str) repetitions_c = np.zeros(len(features_c), dtype=int) pos = 0 for ii in range(len(features)): size = len(features[ii]) groups_c[pos:pos+size] = groups[ii][0] repetitions_c[pos:pos+size] = repetitions[ii] pos += size # Run R with the given template script rscript = importlib_resources.read_text("dclab.lme4", _, script_path = tempfile.mkstemp(prefix="dclab_lme4_", suffix=".R", script_path = pathlib.Path(script_path) rscript = rscript.replace("<MODEL_NAME>", self.model) rscript = rscript.replace("<FEATURES>", arr2str(features_c)) rscript = rscript.replace("<REPETITIONS>", arr2str(repetitions_c)) rscript = rscript.replace("<GROUPS>", arr2str(groups_c)) script_path.write_text(rscript, encoding="utf-8") result = rsetup.run_command((rsetup.get_r_script_path(), script_path)) ret_dict = self.parse_result(result) ret_dict["is differential"] = self.is_differential() ret_dict["feature"] = self.feature ret_dict["r script"] = rscript ret_dict["r output"] = result assert ret_dict["model"] == self.model return ret_dict def get_differential_dataset(self): """Return the differential dataset for channel/reservoir data The most famous use case is differential deformation. The idea is that you cannot tell what the difference in deformation from channel to reservoir, because you never measure the same object in the reservoir and the channel. You usually just have two distributions. Comparing distributions is possible via bootstrapping. And then, instead of running the lme4 analysis with the channel deformation data, it is run with the differential deformation (subtraction of the bootstrapped deformation distributions for channel and reservoir). features = [] groups = [] repetitions = [] # compute differential features for grp in sorted(set([dd[1] for dd in self.data])): # repetitions per groups grp_rep = sorted(set([dd[2] for dd in self.data if dd[1] == grp])) for rep in grp_rep: feat_cha = self.get_feature_data(grp, rep, region="channel") feat_res = self.get_feature_data(grp, rep, region="reservoir") bs_cha, bs_res = bootstrapped_median_distributions(feat_cha, # differential feature features.append(bs_cha - bs_res) return features, groups, repetitions def get_feature_data(self, group, repetition, region="channel"): """Return array containing feature data group: str Measurement group ("control" or "treatment") repetition: int Measurement repetition region: str Either "channel" or "reservoir" fdata: 1d ndarray Feature data (Nans and Infs removed) assert group in ["control", "treatment"] assert isinstance(repetition, numbers.Integral) assert region in ["reservoir", "channel"] for dd in self.data: if dd[1] == group and dd[2] == repetition and dd[3] == region: ds = dd[0] raise ValueError("Dataset for group '{}', repetition".format(group) + " '{}', and region".format(repetition) + " '{}' not found!".format(region)) fdata = ds[self.feature][ds.filter.all] fdata_valid = fdata[~np.logical_or(np.isnan(fdata), np.isinf(fdata))] return fdata_valid def is_differential(self): """Return True if the differential feature is computed for analysis This effectively just checks the regions of the datasets and returns True if any one of the regions is "reservoir". See Also get_differential_features: for an explanation for dd in self.data: if dd[3] == "reservoir": return True return False def parse_result(self, result): resd = result.split("OUTPUT") ret_dict = {} for item in resd: string = item.split("#*#")[0] key, value = string.split(":", 1) key = key.strip() value = value.strip().replace("\n\n", "\n") if key == "fixed effects repetitions": rows = value.split("\n")[1:] reps = [] for row in rows: reps.append([float(vv) for vv in row.split()[1:]]) value = np.array(reps).transpose() elif key == "model converged": value = value == "TRUE" elif value == "NA": value = np.nan value = float(value) except ValueError: ret_dict[key] = value return ret_dict def set_options(self, model=None, feature=None): """Set analysis options""" if model is not None: assert model in ["lmer", "glmer+loglink"] self.model = model if feature is not None: assert dfn.scalar_feature_exists(feature) self.feature = feature def arr2str(a): """Convert an array to a string""" if isinstance(a.dtype.type, np.integer): return ",".join(str(dd) for dd in a.tolist()) elif a.dtype.type == np.str_: return ",".join(f"'{dd}'" for dd in a.tolist()) return ",".join(f"{dd:.16g}" for dd in a.tolist()) def bootstrapped_median_distributions(a, b, bs_iter=1000, rs=117): """Compute the bootstrapped distributions for two arrays. a, b: 1d ndarray of length N Input data bs_iter: int Number of bootstrapping iterations to perform (output size). rs: int Random state seed for random number generator median_dist_a, median_dist_b: 1d arrays of length bs_iter Boostrap distribution of medians for ``a`` and ``b``. See Also From a programmatic point of view, it would have been better to implement this method for just one input array (because of redundant code). However, due to historical reasons (testing and comparability to Shape-Out 1), bootstrapping is done interleaved for the two arrays. # Seed random numbers that are reproducible on different machines prng_object = np.random.RandomState(rs) # Initialize median arrays median_a = np.zeros(bs_iter) median_b = np.zeros(bs_iter) # If this loop is still too slow, we could get rid of it and # do everything with arrays. Depends on whether we will # eventually run into memory problems with array sizes # of y*bs_iter and yR*bs_iter. lena = len(a) lenb = len(b) for q in range(bs_iter): # Compute random indices and draw from a, b draw_a_idx = prng_object.randint(0, lena, lena) median_a[q] = np.median(a[draw_a_idx]) draw_b_idx = prng_object.randint(0, lenb, lenb) median_b[q] = np.median(b[draw_b_idx]) return median_a, median_b
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seminars - Transference theorems, parametric geometry of numbers, and spectra II We define classical exponents of Diophantine approximation attached to real matrices and state various transference theorems between them. We show how the parametric geometry of numbers, introduced by Schmidt and Summerer and further developed by Roy, allows us to give an easy proof of these transference theorems. Furthermore, we explain how this recent theory can be applied to determine the spectra of several exponents of Diophantine approximation (that is, the set of values they take on real entries).
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The largest subsemilattices of the endomorphism monoid of an independence algebra Araújo, João; Bentz, Wolfram; Konieczny, Janusz Linear Algebra and its Applications, 458 (2014), 50-79 http://dx.doi.org/10.1016/j.laa.2014.05.041 (preprint - http://arxiv.org/pdf/1007.4845) An algebra A is said to be an independence algebra if it is a matroid algebra and every map ?:X?A, defined on a basis X of A, can be extended to an endomorphism of A. These algebras are particularly well-behaved generalizations of vector spaces, and hence they naturally appear in several branches of mathematics such as model theory, group theory, and semigroup theory. It is well known that matroid algebras have a well-defined notion of dimension. Let A be any independence algebra of finite dimension n , with at least two elements. Denote by End(A) the monoid of endomorphisms of A. We prove that a largest subsemilattice of End(A) has either 2n?1 elements (if the clone of A does not contain any constant operations) or 2n elements (if the clone of A contains constant operations). As corollaries, we obtain formulas for the size of the largest subsemilattices of: some variants of the monoid of linear operators of a finite-dimensional vector space, the monoid of full transformations on a finite set X, the monoid of partial transformations on X, the monoid of endomorphisms of a free G-set with a finite set of free generators, among others. The paper ends with a relatively large number of problems that might attract the attention of experts in linear algebra, ring theory, extremal combinatorics, group theory, semigroup theory, universal algebraic geometry, and universal algebra.
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Stokes Radius Calculator - Savvy Calculator Stokes Radius Calculator The Stokes radius is a measure of the size of a spherical particle in a fluid, based on its terminal velocity in the fluid. It is particularly useful in fields such as biology, chemistry, and fluid dynamics for understanding particle behavior and interactions. The Stokes radius (rrr) can be calculated using the formula: r=9⋅η⋅v2⋅g⋅(ρp−ρf)r = \sqrt{\frac{9 \cdot \eta \cdot v}{2 \cdot g \cdot (\rho_p – \rho_f)}}r=2⋅g⋅(ρp−ρf)9⋅η⋅v • rrr is the Stokes radius (meters) • η\etaη is the viscosity of the fluid (Pa·s) • vvv is the terminal velocity of the particle (m/s) • ggg is the acceleration due to gravity (9.81 m/s²) • ρp\rho_pρp is the density of the particle (kg/m³) • ρf\rho_fρf is the density of the fluid (kg/m³) How to Use To use the Stokes Radius Calculator: 1. Enter the viscosity of the fluid in Pascal-seconds (Pa·s). 2. Enter the terminal velocity of the particle in meters per second (m/s). 3. Enter the density difference between the particle and the fluid in kilograms per cubic meter (kg/m³). 4. Click the “Calculate” button. 5. The Stokes radius of the particle will be calculated and displayed in meters (m). Suppose we have a particle in a fluid with the following properties: • Viscosity = 0.002 Pa·s • Terminal velocity = 0.1 m/s • Density difference = 50 kg/m³ Using the calculator: 1. Enter 0.002 for viscosity. 2. Enter 0.1 for terminal velocity. 3. Enter 50 for density difference. 4. Click “Calculate.” 5. The Stokes radius will be calculated as approximately 0.018 meters. 1. What is the Stokes radius? □ The Stokes radius is the radius of a spherical particle in a fluid, calculated based on its terminal velocity and fluid properties. 2. Why is the Stokes radius important? □ It helps in determining the size of particles in suspension and understanding their behavior in fluids. 3. Can the Stokes radius be used for non-spherical particles? □ The Stokes radius formula assumes spherical particles. For non-spherical particles, different calculations or models may be required. 4. What are the units of Stokes radius? □ Stokes radius is measured in meters (m). 5. How does density difference affect the Stokes radius? □ A greater density difference between the particle and the fluid results in a larger Stokes radius. 6. Is the Stokes radius the same as hydrodynamic radius? □ No, they are different. Hydrodynamic radius considers the overall size of a particle’s diffusion in a fluid, whereas Stokes radius specifically calculates size based on terminal velocity. 7. Can Stokes radius be used in biological sciences? □ Yes, it is commonly used to determine the size of macromolecules and particles in biological fluids. 8. What is the significance of the terminal velocity in Stokes radius calculation? □ Terminal velocity indicates the speed at which a particle falls through a fluid under gravity, crucial for determining its Stokes radius. 9. Does the calculator account for all factors influencing particle size in fluids? □ The calculator focuses on the basic parameters required for Stokes radius calculation. For specific applications, additional factors may need consideration. 10. Can Stokes radius be experimentally measured? □ Yes, experimental methods involve observing the behavior of particles in fluid flow and relating it to their properties and terminal velocities. The Stokes Radius Calculator provides a straightforward method for determining the size of particles in fluids based on their terminal velocities and fluid properties. By using the formula and inputting the relevant data, users can quickly obtain the Stokes radius, facilitating research and analysis in various scientific disciplines.
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Subtractor | BimStudies.Com A subtractor is a digital circuit that performs subtraction of binary numbers. It’s a fundamental component in digital logic systems and is used to calculate the difference between two binary values. Key components of a subtractor include: Minuend: The number from which subtraction is performed. Subtrahend: The number that is subtracted from the minuend. Borrow (Borrow-in): A signal indicating whether a borrow needs to be taken into account in subtraction. Difference: The result of the subtraction operation. Borrow-out: A signal indicating whether a borrow occurred in the subtraction operation. There are different types of subtractors, including half subtractors, full subtractors, and n-bit parallel subtractors. Half Subractor: A half subtractor is a digital circuit that performs subtraction of two single binary digits (bits). It produces two outputs: a difference (D) and a borrow (B) output. Unlike a full subtractor, a half subtractor doesn’t take into account any borrow from a previous stage, making it suitable for subtracting two individual bits. Block Diagram: Truth Table: Logic Circuit: Full Subtractor: A full subtractor is a digital circuit that performs subtraction of three single binary digits: the minuend bit (A), the subtrahend bit (B), and a borrow-in bit (Bin) from a previous stage. It produces two outputs: a difference (D) and a borrow-out (Bout) output. Unlike a half subtractor, a full subtractor takes into account both the subtrahend and the potential borrow from the previous stage. Block Diagram: Truth Table: Logic Circuit:
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Williams, Llewelyn January 1960 (has links) Thesis (Ph.D.)--Boston University / One of the more powerful and consistent tools available to nature is the phenomenon of alternate freezing and thawing. Mechanically, extraordinary pressures may be involved because of the density differential existing between the liquid and the solid phases of water. Physiologically, there is the availability or non-availability of water to sustain growth. Despite this, catastrophic changes are not to be expected. On the other hand, such a powerful tool must leave its imprint in one manner or another upon the natural landscape. In most arctic and highland areas the imprint is directly discernible. In more moderate climes the imprint is indirectly applied principally as a limiting parameter within an aggregate of generally 1 favorable conditions. The phenomenon of freeze-thaw is a climatic parameter but not a climatic element. Unlike the elements, there is a definite threshold involved; that is, 32 degrees Fahrenheit or 0 degrees Centigrade. At this threshold water may exist in either the liquid or solid state but by the addition or subtraction of heat it can change from one state to the other without a gain or loss in temperature. In the natural environment the terms are not quite so precise. Time for the process to take place, impurities in the water, and the variation of temperature regimes among the many nooks and crannies of the landscape point to the necessity of relaxing the temperature threshold. In this study the zone of 34 degrees F and 28 degrees F is used. Conditions favorable for a thaw are thought to occur when the temperature rises through the zone and conditions favorable for a freeze when the temperature drops through this zone [TRUNCATED] Panturat, Suwanna, January 1987 (has links) 2 Thesis (Ph.D.)--University of Oklahoma, 1987. / Bibliography: leaves 213-219. Kaczmarczyk, Edward Bruce. January 1978 (has links) 3 Thesis (M.S.)--Wisconsin. / Includes bibliographical references (leaves 87-91). Wang, Xiaoxue, 王霄雪 January 2015 (has links) 4 abstract / Mechanical Engineering / Doctoral / Doctor of Philosophy Urban climatology LaJoie, Mark R. 03 1900 (has links) Department of Defense (DoD) climatology products rely mainly on long term means (LTMs) of climate system variables. In this project, we have demonstrated that climatologies based on LTMs can be substantially improved using modern data and methods, especially by accounting for climate variations. We analyzed, and identified mechanisms for, enhanced (suppressed) autumn precipitation in the Horn of Africa (HOA) during El Nino (La Nina) events. El Nino (La Nina) precipitation anomalies were associated with anomalously warm (cool) western Indian Ocean sea surface temperatures, and with 5 anomalously onshore (offshore) moisture transports in the HOA. These transport anomalies supported anomalously strong (weak) precipitation over the HOA. To improve climatological support for DoD operations, we developed and tested a six-step smart climatology process. We applied this process in the context of a notional, unclassified non-combatant evacuation operation (NEO) set in the HOA during autumn of an El Nino year. Using this process, we translated our scientific and operational findings into warfighter impacts. The smart climatology process we have developed is readily adaptable to other regions, seasons, climate variations, and military operations. We have provided a detailed description of our smart climatology process to facilitate its use by DoD agencies. Cannon, Alex Jason 05 1900 (has links) This dissertation develops multivariate statistical models for seasonal forecasting and downscaling of climate variables. In the case of seasonal climate forecasting, where record lengths are typically short and signal-to-noise ratios are low, particularly at long lead-times, forecast models must be robust against noise. To this end, two models are developed. Robust nonlinear canonical correlation analysis, which introduces robust cost functions to an existing model architecture, is outlined in Chapter 2. Nonlinear principal predictor analysis, the nonlinear extension of principal predictor analysis, a linear model of intermediate complexity between multivariate regression and canonical correlation analysis, is developed in Chapter 3. In the case of climate downscaling, the goal is to predict values of weather elements observed at local or regional scales from the synoptic-scale atmospheric circulation, usually for the purpose of generating climate scenarios from Global Climate Models. In this context, models must not only be accurate in terms of traditional model verification statistics, but they must also be able to replicate statistical 6 properties of the historical observations. When downscaling series observed at multiple sites, correctly specifying relationships between sites is of key concern. Three models are developed for multi-site downscaling. Chapter 4 introduces nonlinear analog predictor analysis, a hybrid model that couples a neural network to an analog model. The neural network maps the original predictors to a lower-dimensional space such that predictions from the analog model are improved. Multivariate ridge regression with negative values of the ridge parameters is introduced in Chapter 5 as a means of performing expanded downscaling, which is a linear model that constrains the covariance matrix of model predictions to match that of observations. The expanded Bernoulli-gamma density network, a nonlinear probabilistic extension of expanded downscaling, is introduced in Chapter 6 for multi-site precipitation downscaling. The single-site model is extended by allowing multiple predictands and by adopting the expanded downscaling covariance constraint. Cannon, Alex Jason 05 1900 (has links) This dissertation develops multivariate statistical models for seasonal forecasting and downscaling of climate variables. In the case of seasonal climate forecasting, where record lengths are typically short and signal-to-noise ratios are low, particularly at long lead-times, forecast models must be robust against noise. To this end, two models are developed. Robust nonlinear canonical correlation analysis, which introduces robust cost functions to an existing model architecture, is outlined in Chapter 2. Nonlinear principal predictor analysis, the nonlinear extension of principal predictor analysis, a linear model of intermediate complexity between multivariate regression and canonical correlation analysis, is developed in Chapter 3. In the case of climate downscaling, the goal is to predict values of weather elements observed at local or regional scales from the synoptic-scale atmospheric circulation, usually for the purpose of generating climate scenarios from Global Climate Models. In this context, models must not only be accurate in terms of traditional model verification statistics, but they must also be able to replicate statistical 7 properties of the historical observations. When downscaling series observed at multiple sites, correctly specifying relationships between sites is of key concern. Three models are developed for multi-site downscaling. Chapter 4 introduces nonlinear analog predictor analysis, a hybrid model that couples a neural network to an analog model. The neural network maps the original predictors to a lower-dimensional space such that predictions from the analog model are improved. Multivariate ridge regression with negative values of the ridge parameters is introduced in Chapter 5 as a means of performing expanded downscaling, which is a linear model that constrains the covariance matrix of model predictions to match that of observations. The expanded Bernoulli-gamma density network, a nonlinear probabilistic extension of expanded downscaling, is introduced in Chapter 6 for multi-site precipitation downscaling. The single-site model is extended by allowing multiple predictands and by adopting the expanded downscaling covariance constraint. Trevaskis, Graham Arthur. January 1963 (has links) 8 Thesis (Ed.D.)--Teachers College, Columbia University, 1963. / Typescript; issued also on microfilm. Sponsor: James P. Matthai. Dissertation Committee: Margaret Lindsey. Includes bibliographical references (leaves 206-219). Climatology. Geography Hall, Leonard Frazier, January 1977 (has links) 9 Thesis--Wisconsin. / Vita. Includes bibliographical references (leaves 162-166). Plains Climatology Lettau, Bernhard. January 1961 (has links) 10 Thesis (M.S.)--University of Wisconsin--Madison, 1961. / Typescript. eContent provider-neutral record in process. Description based on print version record. Includes bibliographical references (leaf 19). Bogs. Climatology.
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Algorithm for reducing computational costs in problems of calculation of asymmetrically loaded shells of rotation Algorithm for reducing computational costs in problems of calculation of asymmetrically loaded shells of rotation rotation shells, variable stiffness, asymmetric loading, coefficient prediction in the Fourier method, reduction in computing costs The problem of calculating the shells of rotation of a variable along the meridian of rigidity under asymmetric loading is reduced o a set of systems of one-dimensional boundary value problems with respect to the amplitudes of decomposition of the required functions into trigonometric Fourier series. A method for reducing the number of one-dimensional boundary value problems required to achieve a given accuracy in determining the stress-strain state of the shells of rotation with a variable along the meridian wall thickness under asymmetric load. The idea of the proposed approach is to apply periodic extrapolation (prediction) of the values of the decomposition coefficients of the required functions using the results of calculations of previous coefficients of the corresponding trigonometric series, thus replacing them with some prediction values calculated by simple formulas. To solve this problem, we propose the joint use of Aitken-Steffens extrapolation dependences and Adams method in the form of incremental component, which is quite effective in solving the Cauchy problem for systems of ordinary differential equations and is based on Lagrange and Newton extrapolation dependences. The validity of the proposed approach was verified b the results of a systematic numerical experiment by predicting the values of the expansion coefficients in the Fourier series of known functions of one variable. The approach is quite effective in the calculation of asymmetrically loaded shells of rotation with variable along the meridian thickness, when the coefficients of decomposition of the required functions into Fourier series are functions of the longitudina lcoordinate and are calculated by solving the corresponding boundary value problem. In this case, the approach allows solving solutions of differential equations for the amplitudes of decompositionin to trigonometric series only for individual "reference" harmonics, and the amplitudes for every third harmonic can be calculated by interpolating their values for all node integration points of the corresponding boundary value problem. This significantly reduces the computational cost of obtaining the solution as a whole. As an example, the results of the calculation of the stress-strain state of a steel annular plate under asymmetric transverse loading are given. Biderman V. L. Mechanics of thin-walled structures / V. L. Biderman. – Mashinostroenie. –Мosscow, 1977. – 488 p. (in Russian). Grigorenko Ya. M. Methods for calculating shells. Theory of shells of variable rigidity. Vol. 4 / Ya. M. Grigorenko, A. T. Vasilenko. – Naukova Dunka, Kiev, 1981. – 544 p. (in Russian). Grigorenko A. Ya. Stress-strain state of shallow rectangular shells of variable thickness under various boundary conditions / A. Ya Grigorenko, N. P Yaremchenko., C. N. Yaremchenko // Bul. of NAS Ukraine. – Vol. 6. – 2016. – P. 31-37 (in Ukrainian). Dzyuba A. P. Calculation algorithm on the basis of a discrete-continuous approach for cylindrical shell of variable rigidity in circular direction / A. P. Dzyuba, I. A. Safronova, L. D. Levitina // Problems of computational mechanics and strength of structures, Collection of scientific articles. –Vol. 30. – 2019. – P. 53-67 (in Ukrainian). Myachenkov V. I. Calculation of composite shell structures on a computer: Reference / V. I. Myachenkov I. V. Grigoryev. – Mashinostroenie. – Moscow, 1981. – 212 p. (in Russian). Sineva N. F. Calculation of a cylindrical shell of variable stiffness interacting with a nonlinear elastic base / N. F. Sineva, F. S. Selivanov, D. V. Nikityuk // Bull. Saratov State Techn. Un-ty. Ser.: Construction and architecture. – Vol. 4(60). – Iss 2. – 2011. – P. 15-21 (in Russian). Models and algorithms for optimization of elements of nonuniform shell structures, in N. V. Pоlyakov (Eds.) / A. P. Dzyuba, V. N. Sirenko, A. A. Dzyuba and I. A. Safronova // Actual problems of mechanics: Monograph, Lira, Dnipro, 2018. – P. 225-243 (in Ukrainian). Ovchinnikov I. G. Thin-walled structures under conditions of corrosion wear: Calculation and optimization / I. G. Ovchinnikov, Yu. M. Pochtman. – DNU. – Dnepropetrovsk, 1995. – 190 p. (in Russian). Dashchenko A. F. ANSYS in the problems of mechanical engineering: monograph. second ed. / A. F. Dashchenko, D. V. Lazareva, N. G. Suryaninov. – BURUN and K0. – Kharkiv, 2011. – 504 p. (in Russian). Alyamovsky A. A. SolidWorks/COSMOSWorks. Finite Element Engineering Analysis. / A. A. Alyamovsky. – DMK Press. – Moscow, 2004. 432 p. (in Russian). Grigorenko Ya. M. Solving the problems of shell theory on a computer / Ya. M. Grigorenko, A. P. Mukoed. – High School. – Kiev, 1979. – 280 p. (in Russian). Mossakovsky V. I. Contact Interactions of Elements of Shell Structures / V. I. Mossakovsky, V. S. Hudramovich, E. M. Makeev. – Naukova Dunka. – Kiev, 1988. – 288 p. (in Russian). Emel’yanov I. G. Application of discrete Fourier series to the stress analysis of shell structures / I. G. Emel’yanov // Computational Continuum Mechanics. – Vol 8(3). – 2015. – P. 245-253. (in Hudramovich V. S. Contact interaction and optimization of locally loaded shell structures / V. S. Hudramovich, A. P. Dzyuba // Journal of mathematical Science - Springer Science + Business media. – 2009. – P. 231-245. Strength. Sustainability. Fluctuations: Handbook, vol.1. / (Eds.) I. A. Birger, Ya. G. Panovko. – Mashinostroenie. – Moscow, 1968. – 821 p. (in Russian). Tolstov G. P. Fourier Series / G. P. Tolstov. – Fizmatgiz. – Moscow, 1960. – 392 p. (in Russian). Godunov S. K. On the numerical solution of boundary value problems for systems of linear ordinary differential equations / S. K. Godunov // Advances in Mathematical Sciences. – Vol. 16. – Iss.3 (99). –1961. – P. 171-174 (in Russian). Bulakajev P. I. An algorithm for the prediction of search trajectory in nonlinear programming problems optimum design / P. I. Bulakajev, A. P. Dzjuba // Structural Optimization: Research Journal of Intern. Society of Struct. and Multidisciplinary Optimiz. – Springer – Verlag. –Vol. 13(2, 3). – 1997. –P. 199-202. Dzyuba A. P. An algorithm for reducing the computational cost of using the Fourier method in the problems of shell structural mechanics / A. P. Dzyuba, O. O. Bobilev, P. I. Bulakaev // Bull. of Dnepropetrovsk state university. Vol. 2(2). – 1999. – P. 47-57 (in Russian). Shamansky V. E. Methods of numerical solution of boundary value problems on a computer / V. E. Shamansky. – Academia Science USSR. – Kiev. – Part. 1, 1963. – 196 p. – Part. 2, 1966. – 242 p. (in This work is licensed under a Creative Commons Attribution 4.0 International License. Authors retain copyright and grant the journal right of first publication with the work simultaneously licensed under a Creative Commons Attribution License that allows others to share the work with an acknowledgement of the work's authorship and initial publication in this journal.
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Flip a Column of Data Vertically - Free Excel Tutorial This post will guide you how to flip a column of data vertically preserving the original formatting and formulas in Excel. How do I reverse a column of data order in a column arranged alphabetically or from largest to smallest. How to use the Excel Sort feature to reverse the data order? Or how to use a VBA macro code to flip a column of data vertically in Excel. If you want to flip or reverse a column of data that is already in numeric or alphabetical order, and you can easily to flip it with Sort feature. And if the column data is not sorted and you can try to create a helper column to achieve the result. 1. Flip a Column of Data with Sort Feature #1 select the column that you want to flip, right click on it, and select Insert to add helper column. Choose Entire column radio button. Then insert 1,2,3,4,…in helper column. #2 select your helper column, then go to the Data tab, click the Sort Largest to Smallest command under Sort&Filter group. #3 you will see that the numbers in help column will be sorted and the data of column also be reversed. Now you can delete the help column. Or you can also hide the help column in case you want to restore the data of column. #4 let’s see the result 2. Flip a Column of Data Vertically with Function You can also create an excel formula based on the INDEX function and the ROWS function to flip a column of data vertically. Type the following formula in the formula box of cell C1, then press enter Let’s see the result as below: 3. Flip a Column of Data Vertically with VBA Macro You can also write an Excel VBA Macro to flip a column of data vertically quickly. Just do the following steps: #1 click on “Visual Basic” command under DEVELOPER Tab. #2 then the “Visual Basic Editor” window will appear. #3 click “Insert” ->”Module” to create a new module. #4 paste the below VBA code into the code window. Then clicking “Save” button. Sub FlipCol() Dim R As Range Dim W As Range Dim A As Variant Dim i As Integer Dim j As Integer Dim k As Integer On Error Resume Next titleS = "Flip data in a column" Set W = Application.Selection Set W = Application.InputBox("select a range that you want to reverse", titleS, W.Address, Type:=8) A = W.Formula For j = 1 To UBound(A, 2) k = UBound(A, 1) For i = 1 To UBound(A, 1) / 2 xTemp = A(i, j) A(i, j) = A(k, j) A(k, j) = xTemp k = k - 1 W.Formula = A End Sub #5 back to the current worksheet, then run the above excel macro. Click Run button. #6 select a range that you want to flip. #7 click OK button, then check the result. 4. Video: Flip a Column of Data Vertically This Excel video tutorial, we’ll explore three methods to flip a column of data vertically. We’ll start by using a formula with the INDEX and ROWS functions, followed by leveraging the Sort feature, and finally, we’ll delve into a VBA Macro. 5. Related Functions • Excel INDEX function The Excel INDEX function returns a value from a table based on the index (row number and column number)The INDEX function is a build-in function in Microsoft Excel and it is categorized as a Lookup and Reference Function.The syntax of the INDEX function is as below:= INDEX (array, row_num,[column_num])… • Excel ROWS function The Excel ROWS function returns the number of rows in a cell reference.The syntax of the ROWS function is as below:= ROWS(array)…
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For homework submissions, I require that all work be completed - Essay Writing For homework submissions, I require that all work be completed in excel, which means that you will use the excel functions to arrive at the answer. I should be able to click on the cell that holds your answer, and the formula to arrive at this answer will be shown in the function ribbon. I will not accept answers that were simply re-typed into the excel workbook. Chapter 1. Year Stock T Stock B 1 0.19 0.08 2 0.08 0.03 3 -0.12 -0.09 4 -0.03 0.02 5 0.15 0.04 a. Compute the arithmetic mean annual rate of return for each stock. Which stock is most desirable by this measure? b. Compute the standard deviation of the annual rate of return for each stock. (Use Chapter 1 Appendix if necessary.) By this measure, which is the preferable stock? c. Compute the coefficient of variation for each stock. (Use the Chapter 1 Appendix if necessary.) By this relative measure of risk, which stock is preferable? d. Compute the geometric mean rate of return for each stock. Discuss the difference between the arithmetic mean return and the geometric mean return for each stock. Discuss the differences in the mean returns relative to the standard deviation of the return for each stock. 7. A stockbroker calls you and suggests that you invest in the Lauren Computer Company. After analyzing the firm’s annual report and other material, you believe that the distribution of expected rates of return is as follows: Possible Rate of Return Probability -0.60 0.05 -0.30 0.20 -0.10 0.10 0.20 0.30 0.40 0.20 0.80 0.15 Compute the expected return [E(Ri)] on Lauren Computer stock. 9. During the past year, you had a portfolio that contained U.S. government T-bills, long-term government bonds, and common stocks. The rates of return on each of them were as follows: U.S. government T-bills 5.50% U.S. government long-term bonds 7.50 U.S. common stocks 11.60 During the year, the consumer price index, which measures the rate of inflation, went from 160 to 172 (1982 – 1984 = 100). Compute the rate of inflation during this year. Compute the real rates of return on each of the investments in your portfolio based on the inflation rate. 12. Assume that the consensus required rate of return on common stocks is 14 percent. In addition, you read in Fortune that the expected rate of inflation is 5 percent and the estimated long-term real growth rate of the economy is 3 percent. What interest rate would you expect on U.S. government T-bills? What is the approximate risk premium for common stocks implied by these data? Chapter 2. 4. a. Someone in the 36 percent tax bracket can earn 9 percent annually on her investments in a tax-exempt IRA account. What will be the value of a one-time $10,000 investment in 5 years? 10 years? 20 years? b. Suppose the preceding 9 percent return is taxable rather than tax-deferred and the taxes are paid annually. What will be the after-tax value of her $10,000 investment after 5, 10, and 20 years? 5. a. Someone in the 15 percent tax bracket can earn 10 percent on his investments in a tax-exempt IRA account. What will be the value of a $10,000 investment in 5 years? 10 years? 20 years? b. Suppose the preceding 10 percent return is taxable rather than tax-deferred. What will be the after-tax value of his $10,000 investment after 5, 10, and 20 years? 6. Assume that the rate of inflation during all these periods was 3 percent a year. Compute the real value of the two tax-deferred portfolios in problems 4a and 5a.
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Speed Conversion Speed Conversion Conversion between nearly 75 different Metric, English, Imperial, Indian Speed (and Velocity) measurement units. Supports conversions between Kmph, Mph, Knot, Mach, Benz, Yard per hour, Foot per second, Meter per second, Megameter per hour, Speed of Light, Speed of Sound and many more Metric, English and Imperial units. Example: , , , , , , , , , , , and many more. Acceleration Conversion Conversion between nearly 25 different Acceleration measurement units. Supports conversion between Centigal, Decigal, Milligal, G-Unit (G), Gal, Galileo, Gn, Grav, Meter per second-square, Centimeter per second-square, Inch per second-square, Decimeter per second-square, Dekameter per second-square, Foot per second-square, Hectometer per second-square, Kilometer per hour-second, Kilometer per second-square, Mile per hour-minute, Mile per hour-second, Mile per second-square, Millimeter per second-square and many more Torque Conversion Conversion between nearly 20 different Torque measurement units. Supports conversion between Dyne Centimeter, Gram Centimeter, Kilogram Centimeter, Kilogram Meter, Kilonewton Meter, Kilopond Meter, Meganewton Meter, Micronewton Meter, Millinewton Meter, Newton Centimeter, Newton Meter, Ounce Foot, Ounce Inch, Pound Foot, Pound Inch, Poundal Foot and many more Dynamic Viscosity Conversion Conversion between nearly 25 different Dynamic Viscosity measurement units. Supports conversion between Poise, Decipoise, Centipoise, Millipoise, Poiseuille (France), Reyn, Newton second/m², Pascal second, Millipascal second, Dyne second/cm², Gram-Force second/cm², Gram/centimeter second, Kilogram-Force second/m², Kilogram/meter hour, Kilogram/meter second, Pound-Force second/f², Pound-Force second/in², Pound/foot hour, Pound/foot second, Poundal hour/feet², Poundal second/feet², Slug/foot second and many more Kinematic Viscosity Conversion Conversion between nearly 30 different Kinematic Viscosity measurement units. Supports conversion between Stokes, Centistokes, Lentor, Liter per Centimeter-Day, Liter per Centimeter-Hour, Liter per Centimeter-Minute, Liter per Centimeter-Second, Poise Cubic Centimeter per Gram, Square-Centimeter per Day, Square-Centimeter per Hour, Square-Centimeter per Minute, Square-Centimeter per Second, Square-Feet per Day, Square-Feet per Hour, Square-Feet per Minute, Square-Feet per Second, Square-Inch per Day, Square-Inch per Hour, Square-Inch per Minute, Square-Inch per Second, Square-Meter per Day, Square-Meter per Hour, Square-Meter per Minute, Square-Meter per Second, Square-Millimeter per Day, Square-Millimeter per Hour, Square-Millimeter per Minute, Square-Millimeter per Second and many more Viscosity of Oil, Water and Glycerin Conversion Conversion of viscosity at different temperature of Oil, Water, Glycerin and thick fluids between nearly 15 different measurement units. Supports conversion between Poise, Centipoise, Water at 20C, Water at 40C, Heavy Oil at 20C, Heavy Oil at 40C, Glycerin at 20C, Glycerin at 40C, SAE 5W at -18C, SAE 10W at -18C, SAE 20W at -18C, SAE 5W at 99C, SAE 10W at 99C, SAE 20W at 99C and many Angular Velocity Conversion Conversion between degree/day, degree/hour, degree/minute, degree/second, radian/second, radian/day, radian/hour, radian/minute, revolution/day, revolution/hour, revolution/minute, revolution/ second and more Angular Acceleration Conversion Conversion between radian/second^2, radian/minute^2, degree/second^2, degree/minute^2, revolution/second^2, revolution/minute^2 Moment of Inertia Conversion Conversion between nearly 20 different Moment of Inertia measurement units. Support conversion between Kg-m^2, Kg-cm^2, Kg-mm^2, g-cm^2, g-mm^2, oz-in^2, lb-ft^2, lb-in^2, slug-ft^2 and many more Speed of common Devices, Peripheral, Protocols Find speed of sound, light, electricity, speed of different computer devices, peripherals, protocols, upload and download speed, bit rate and more Escape Velocity of different Planets, Moons, Comets and Asteroids Find escape velocity, surface gravity, mass of different solar system objects, planets, moons, comets, asteroids, centaurs, kuiper belt objects and more. Example: Find Escape Velocity of , , , , , , , , , , , , , , , , , , , , , , , , Speed of common Animals Find top speed of different animals, top speed of common land animals, top speed of common birds, top speed of aquatic (sea) animals, top speed of different insects, fastest animal, fastest mammal, fastest fish, fastest bird and many more. Also see ...
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Integrability and complexity in quantum SciPost Submission Page Integrability and complexity in quantum spin chains by Ben Craps, Marine De Clerck, Oleg Evnin, Philip Hacker This is not the latest submitted version. This Submission thread is now published as Submission summary Authors (as registered SciPost users): Marine De Clerck Submission information Preprint Link: scipost_202307_00036v1 (pdf) Code repository: https://doi.org/10.5281/zenodo.7876467 Data repository: https://doi.org/10.5281/zenodo.7876467 Date submitted: 2023-07-28 14:48 Submitted by: De Clerck, Marine Submitted to: SciPost Physics Ontological classification Academic field: Physics Specialties: • Mathematical Physics • Quantum Physics There is a widespread perception that dynamical evolution of integrable systems should be simpler in a quantifiable sense than the evolution of generic systems, though demonstrating this relation between integrability and reduced complexity in practice has remained elusive. We provide a connection of this sort by constructing a specific matrix in terms of the eigenvectors of a given quantum Hamiltonian. The null eigenvalues of this matrix are in one-to-one correspondence with conserved quantities that have simple locality properties (a hallmark of integrability). The typical magnitude of the eigenvalues, on the other hand, controls an explicit bound on Nielsen's complexity of the quantum evolution operator, defined in terms of the same locality specifications. We demonstrate how this connection works in a few concrete examples of quantum spin chains that possess diverse arrays of highly structured conservation laws mandated by integrability. Current status: Has been resubmitted Reports on this Submission Report #2 by Anonymous (Referee 2) on 2023-11-13 (Invited Report) • Cite as: Anonymous, Report on arXiv:scipost_202307_00036v1, delivered 2023-11-13, doi: 10.21468/SciPost.Report.8099 1 - tackles an ambitious and timely problem: complexity of chaotic vs. integrable dynamics 2 - focuses on a computable upper bound to complexity, which exhibits different behavior in chaotic/integrable models 3 - very detailed study of a very interesting object -the 'Q-matrix'- which allows to probe the local conservation laws and has implications on the complexity bound 4 - analytical results on the Q-matrix from random matrix 5 - extensive numerical checks in non-trivial interacting integrable spin chains 6 - very carefully written 1 - complete lack of conciseness 2 - if one reads only this paper, then it is not quite clear what is new or not new here, compared to [SciPost Phys. 13, 090 (2022)] by the same authors Motivated by the notion of complexity of unitary evolutions, the authors study an upper bound on Nielsen's complexity -Eq. (2.24)- related to a matrix $Q$ of size $\dim(\mathcal{H}) \times \dim(\ mathcal{H})$ defined for a subspace of 'easy' operators (such as local or few-body operators). The spectrum of the matrix $Q$ is claimed to be highly sensitive to integrability, in particular its kernel is directly tied to the existence of local conservation laws. The authors make general conjectures about properties of $Q$, which they relate to the behavior of their upper bound on complexity, especially to the plateau that the bound displays at long time. The main claims are supported by analytical calculations for GUE random matrices, in the case of chaotic dynamics (Section 3). Extensive numerical checks are then provided both for chaotic and integrable spin chains (Section 4). I find that the results are very substantial and exciting. They open a new pathway towards the fundamental goal of characterizing the impact of conservation laws on the complexity of quantum many-body dynamics. I also think they provide a synergetic link between research on complexity of quantum evolution operators, which so far has largely focused on chaotic models and/or fully connected models like the SYK model, and interacting integrable spin chains. In my opinion, the only problem of this paper is the fact that it is too long. It took me a very long time to go through it, and while I feel that it was very instructive and inspiring, I am also a bit frustrated because the reading could have been faster and less painful if the manuscript were organized more clearly. Several sections are extremely long, with no clear substructure. For instance, Section 2.1 is 10 pages long, presented as a whole block. It starts by discussing the definition of Nielsen's complexity, then some of its properties and equivalence/differences with other notions of complexity, then its geometric interpretation, then discusses a warm-up calculation of complexity in a simple case, then some generic properties of its short-time behavior, then its long-time behavior and plateau, in relation with typical distances to hybercube lattices, then refers to Ref. [6] for some observation about chaotic vs. integrable dynamics, then comes back to the general discussion of the problem with a penalty factor, then drops Eq. (2.24) -which plays a central role in the whole paper-, then discusses again distances in hypercubic lattices, then makes a long digression about some known numerical methods to solve the 'CVP', then repeats some information about distances in hypercubic lattices, then briefly mentions with Eq. (2.28) that the kernel of the matrix Q encodes the conservation laws of the model -a central point for the rest of the paper-, then discusses whether or not it is a good choice to declare that the identity operator is local, and then provides a loose discussion of integrability, repeating information that already appeared above, and finally briefly hints at the results of Sections 3 and 4. Please... It would not be difficult to break this long section (and other similar sections, such as 4.1.3) that contain digressions and repetitions, into smaller, more focused, subsections or paragraphs, each with a clear title. In addition, it would help the reader to have a clear distinction between the results of Ref. [6] that are partially summarized in section 2.1, and the new results specific to this new paper. In summary, I am happy to recomment publication of these results in Scipost Physics, but I would encourage the authors to try to make the manuscript a bit easier to read, if possible. Requested changes See above. If the authors could give more structure to the manuscript (subsections, paragraph titles, etc), and perhaps also trim it in order to avoid unnecessary repetitions, I think it could really make the manuscript more accessible. Otherwise, the manuscript is very carefully written. I caught only a few typos: - before Eq. (3.10), it seems that a word is missing: perhaps 'We simplify the computation by working to leading order in $1/D$ and we define' - before Eq. (4.3): 'these non-trivial site' -> 'sites' - after Eq. (4.14): 'The former has an explicit U(1) symmetry' sounds weird, because both the former and the latter have that symmetry here Finally, in the captions of the figures, it is said many times that what is shown is 'the complexity' (see e.g. caption of Fig. 6). It would be clearer to recall that it is the 'complexity bound' that is shown Report #1 by Anonymous (Referee 1) on 2023-9-29 (Invited Report) • Cite as: Anonymous, Report on arXiv:scipost_202307_00036v1, delivered 2023-09-29, doi: 10.21468/SciPost.Report.7876 1. Important and relevant results on the connection between complexity and dynamics. 2. Clearly written. Also well contrasted with previous results in the literature (especially Ref. [6] by the same authors) 3. Many practical examples for integrable and chaotic dynamics. This manuscript proposes a new interesting technique to distinguish between chaotic and integrable models based on the concept of complexity. The topic of complexity in many body quantum systems is very timely and of fundamental importance. The same is also the case for the quest of setting apart chaotic and integrable dynamics by means of simple and computable quantities. It is very remarkable to establish a connection between these two subjects. The manuscript is extremely well organized with a clear introduction where the main results are summarized. Sec. 2 contains the definitions of the quantities of interest. In Secs. 3 and 4 the main results for chaotic and integrable dynamics are derived. Many examples are reported to corroborate the results and ideas. I have no doubts concerning the publication of this interesting manuscript in Scipost Physcis. However, I would like that before the authors consider to comment on the point raised below. In the literature, several other quantities to distinguish chaotic and integrable dynamics have been introduced and discussed. Among them, a very important and effective one is the local operator entanglement. Indeed it has been shown (see V. Alba, J. Dubail, and M. Medenjak, Phys. Rev. Lett. 122, 250603 (2019) and references therein) that such quantity grows at most logarithmically in time for integrable models while it grows linearly in chaotic ones. The authors should mention this fact and comment about possible connections between their results and the local operator entanglement.
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Power Measurements on AC-DC Power Supplies AC to DC power supplies are fully integrated into our daily lives. They are the heart of all our electronics, providing energy for the everyday devices that we rely on. A cell phone charger is a power supply, but so are the power electronic circuits embedded inside almost any electronics or appliance connected to the AC line. This near universal use of power supplies means that although the power that each individual supply draws is small, the cumulative effects are very significant. For this reason it is desirable, and often required by regulation, to control power supply characteristics such as efficiency. AC –DC Power supplies convert electrical energy from their AC line input to provide DC outputs that are: • Isolated from the dangerous high-voltage, high-capacity AC line. • Smoothed and low-noise DC voltage. • Regulated to be largely independent of input voltage changes. • Current controlled to avoid damage to the load (especially batteries) and to the power supply itself Typically, an AC to DC power supply converts the AC line (110 220V 50/60Hz) into low voltage (12, 5, 3V) DC. Power supply designers strive to improve the efficiency of their designs while maintaining specified performance over a range of input and load conditions and complying with demanding international regulations for harmonic current content, standby power, safety and EMC. This application note is intended to assist engineers who design and test power supplies and other equipment connected to the AC line make power related measurements accurately, quickly and safely. 2. Power Supply Measurements Essential Background AC Power is the product of RMS input voltage, RMS input current and power factor. The power supply input voltage is the normal AC line in most cases and its shape is near sine-wave with low distortion. The input current may be highly distorted and rich in harmonic content. Most power supplies have input circuits consisting of arectifier followed by a smoothing capacitor to produce a high voltage DC. The capacitor maintains the DC voltage when the AC line input voltage is falling and is charged only when the AC voltage is higher than the DC ‘bus’ voltage. The AC input current is thus no sine wave, but only flows near the peak of the voltage sinewave. An example of a typical switch mode power supply is shown in Figure 1. Figure 2 shows typical input waveforms of a power supply. Figure 1. Switch-mode power supply simplified schematic. Figure 2. Vrms and Arms waveform. This distorted current increases the RMS current drawn from the AC line, but does no useful work and does not contribute to the real power (watts) drawn or used by the power supply. This extra, wasted current, multiplied up by the millions of power supplies in use means that the the AC supply and distribution system must be larger than required to supply the necessary power (watts) and that power is wasted providing the excess current by resistive losses in the transmission and distribution system. This is one reason why the current distortion drawn by power supplies is limited by international regulations such as IEC61000-3-2. Most modern power supply designs mitigate this problem by employing some form of wave shaping circuit on their input, often called a (PFC) Power Factor Correction circuit. An example of PFC implementation is shown in Figure 3. These circuits help shape the input current such that the input impedance of the power supply appears fairly linear to the power line, similar to a resistive load. By doing so, the problems of distorted input current waveforms can be significantly improved. Another key issue, most designers face today is measuring and certifying energy consumption and efficiency on their class of device to comply with institutions like ENERGY STAR™ and CEC. Figure 3. Simplified view of an SMPS power supply (primary side only) and its power quality measurement test points. Simultaneous input VAC and IAC readings are necessary for power quality In order to satisfy these requirements for low input current distortion and harmonics, high efficiency and low average and standby power consumption its important to understand how to measure them in accordance with internationally accepted practice. Figure 4. Distorted input current sampling. 3. Critical Power Supply Measurements It's important for a power supply designer to keep check of all these parameters that are so crucial to a power supply’s performance proper operation and compliance to specification. Some of the most important power measurements a designer has to make are listed as follows: 4. Making Power Supply Input Measurements Using a Tektronix – PA1000 Tektronix PA1000 is a single channel high precision Power Analyzer that helps makes complex power measurements on a power supply very easy. This section will discuss a practical example of making power supply input measurements on a typical power supply. The DUT (Device under test) for this example is a generic laptop charger. The standard current inputs of a power analyzer will measure a large range of current, from milli-amps to 20 or 30 amps RMS. This is suitable for most power supplies up to 3kW. A single power analyzer wattmeter input channel consists of a voltage input pair (V[HI] and V[LO]) and a current input pair (A[HI] and A[LO]). These connections are simplified by use of a break-out box that makes the analyzer connections with 4mm safety connectors and provides a standard AC outlet for connection to the power supply. Make the connections as shown in the Figure 5. Figure 5. PA1000 wiring diagram. Connect the PA1000 and breakout box using the color coded safety-banana lead set as seen in picture For this example a generic laptop charger was plugged in the AC socket on the breakout box and allowed to charge a laptop PA1000 will start displaying results as soon as everything is plugged in and powered on as shown in Figure 6 Figure 6. Result screen. Figure 7 To view more parameters press the "Menu" button on the front panel and navigate to "Measurements" as shown in Figure 8 Figure 8. Measurements menu. Options available include crest factors, power factor, current and voltage harmonics, THD, active and reactive power and many more. Please refer to the user manual for the list of all available Figure 9 shows an example of some of these measurements on the display. Figure 9. More measurement choices. Voltage, current and power waveforms and harmonic bar charts can also be viewed as seen in following Figure 10 and Figure 11 respectively by pressing the Graph button or accessing the Graph Menu. Figure 10. Graph display Figure 11. Harmonic bar chart display. One of the important features that most of the power supply designers and testing labs need is ability to log data efficiently over extended periods of time. PA1000 provides a standard USB port on the front panel for data logging which can be very handy for performing extended testing. The test results can be exported as an excel file and used for presentation or analysis as shown in the Figure 12. Figure 12. Data logged to Excel file. The best and the most convenient way to test power supply is to use the included PWRVIEW software. 5. PWRVIEW Software PWRVIEW software compliments Tektronix Power Analyzers and is available as free download on Tektronix website. PWRVIEW offers easy, wizard-driven test solutions for power supply, standby power and many other target applications For download and installation info go to www.tek.com/power-analyzer-series/pa1000 Tektronix Power Analyzers remote operation Transfer, view, analyze, record and export measurement data in real-time from the instrument, including waveforms, harmonic bar charts, and stand power trend charts. Default applications and standard compliance tests Figure 13 shows the screen shot of PWRVIEW PC software Figure 13. PWRVIEW software. PWRVIEW makes is very easy to monitor, record and analyze the critical power measurement that have been discussed earlier The following section will demonstrate how to make efficiency measurements using two PA1000 using PWRVIEW software. 6. Power Supply Efficiency Measurements Two wattmeter channels are required for performing the efficiency measurements. Either a multichannel Power Analyzer like PA4000 or two PA1000 connected with PWRVIEW can do the job. Two PA1000’s are used for this example to demonstrate efficiency measurements on a power supply. Connect two PA1000’s, one on input and one on output of the power supply under test making connections as demonstrated earlier. Connect both PA’s to the computer via USB, Ethernet or GPIB and add them to PWRVIEW. Once added, PWRVIEW will display two tabs for two Power Analyzers connected as shown in Figure 14. Figure 14. Adding two PA1000 simultaneously For calculating efficiency, just click on the Configure radio tab on the top ribbon. This will open a new tab for input and output selection as shown in Figure 15. Figure 15. Efficiency configuration. Once proper groups are selected for input and output just click on the radio button – Enable Watts Efficiency Measurements. Once done, click on the Measure tab on the top and then click start on the Measure page. The Measurement grid will start making efficiency measurement based on the input and output wattage as seen in the Figure 16. Figure 16. Measurement grid with efficiency. PWRVIEW can also be used to view waveforms and bar charts as shown in following Figure 17 and Figure 18 respectively. Figure 17. Waveform display. Figure 18. Harmonic bar chart display. 7. Standby Power Compliance Tests The power consumed by power supply while they are not loaded is called standby power. The standby power is trivial when considered for individual power supplies but when equated as a whole over a longer period of time and considering the number of power supplies plugged in to sockets all the time, becomes a significant waste of energy Many programs are already in place to reduce standby power including ENERGY STAR and the EU Eco Directive. The scope of these programs continues to grow and the level of standby power in Watts necessary to achieve compliance steadily falls. PWRVIEW software offers engineers a straight forward and accurate tool to conduct a full-compliance standby power test and get to the market faster and cheaper. Figure 19 shows a snapshot of a full compliance standby test in action. Figure 19. Full compliance standby power test. Figure 20. PWRVIEW PC software charts harmonics and compares to limits. For more details on standby power and full compliance stand by test, please refer to Standby Power app note on Harmonics Limits Using PC software coupled to the power analyzer, harmonics measurements may be quickly and conveniently recorded and compared to the limits of IEC61000-3-2 and others. Software features such as PDF report export provide complete reporting functions for power supply conformance measurements. 8. Definitions 8.1. RMS (Root Mean Squared Value) The RMS value is the most commonly used and useful means of specifying the value of both AC voltage and current. The RMS value of an AC waveform indicates the level of power that is available from that waveform and is equivalent DC at the same voltage. This is one of the most important attributes of any AC source. The calculation of an RMS value can best be described by considering an AC current waveform and its associated heating effect such as that shown in Figure 21. Figure 21. Root Mean Square waveform. The current (Amp) is considered to be flowing through a resistance; the heating effect at any instant is given by the equation: By dividing the current cycle into equally spaced coordinates (samples), the variation of the heating effect with time can be determined as shown in Figure 2 above. The average heating effect (power) is given by: To determine the equivalent value of current that would produce the average heating effect value shown above, then the following applies: The RMS value of the current = the square root of the mean of the squares of the current. This value is often termed the effective value of the AC waveform, as it is equivalent to the direct current that produces the same heating effect (power) in the resistive load. It is worth noting that for a sinusoidal waveform: 8.2 Average Value The average value of a waveform such as that shown in Figure 22 is given by: It is clear the average value can only have real meaning over one half cycle of the Waveform, for a symmetrical waveform, the mean or average value over a complete cycle is zero. Most simple multi-meters determine AC values by full-wave rectification of the AC waveform, followed by a calculation of the mean value. Such meters; however, will be calibrated in RMS and will make use of the known relationship between RMS and average for a sinusoidal waveform, i.e. However, for waveforms other than a pure sine wave, the readings from such meters will be invalid. RMS values will need to be calculated using techniques shown in Figure 22. Figure 22. Average calculated methodology. 8.3 Real and Apparent Power (W & VA) A sinusoidal voltage source of, say, 100V RMS is connected to a resistive load of, say, 100Ω, then the voltage and current can be depicted as in Figure 23 and are said to be "in phase". The power that flows from the supply to the load at any instant is given by the value of the product of the voltage and the current at that instant, as illustrated in Figure 23. Figure 23. Voltage and current phase waveform. From this, it can be seen that the power flowing into the load fluctuates (at twice the supply frequency) between 0 and 200W and that the average power delivered to the load equals 100W which is what one might expect from 100V RMS and a resistance of 100 Ω. However, if the load is reactive (i.e. contains inductance or capacitance as well as resistance) with an impedance of 100 Ω, then the current that flows will still be 1A RMS but will no longer be in-phase with the voltage. This is shown in Figure 24 for an inductive load with the current lagging by 60 degrees. Although the power flow continues to fluctuate at twice the supply frequency, it now flows from the supply to the load during only a part of each half cycle—during the remaining part, it actually flows from the load to the supply. The average net flow into the load, therefore, is much smaller than in the case of a resistive load as shown in Figure 23— with only 50W of useful powering delivered into the inductive load. Figure 24. The power that flows from the supply to the load. In both of the above cases the RMS voltage was equal to 100V RMS and the current was 1A RMS. The product of these two values is the apparent power delivered into the load and is measured in VA as The real power delivered has been shown to depend on the nature of the load. It is not possible to determine the value of real power from the knowledge of RMS voltage and the use of a true AC power meter, capable of computing the product of the instantaneous voltage and current values and displaying the mean of the result is required. VA is often measured to ensure that the ac supply has sufficient capacity. Figure 24. Power factor waveform. 8.4 Power Factor It is clear that, in comparison with DC systems, the transferred AC power is not simply the product of the voltage and current values. A further element, known as the power factor must also be taken into consideration. In the previous example (real and apparent power) with an inductive load, the power factor is 0.5 because the useful power is exactly one half of the apparent power. We can therefore define power factor as: In the case of sinusoidal voltage and current waveforms, the power factor is actually equal to the cosine of the phase angle (q) between the voltage and current waveforms. For example, with the inductive load described earlier, the current lags the voltage by 60 degrees. It is for this reason that power factor is often referred to as cosq. However, it is important to remember that this is only the case when both voltage and current are sinusoidal [Figure 24 (I1 and I2)] and that power factor is not equal to cosq in any other case [Figure 24 (I3)]. This must be remembered when using a power factor meter that reads cosq, as the reading will not be valid except for pure sinusoidal voltage and current waveforms. A true power factor meter will compute the ratio of real to apparent power as described in section discussing real and apparent power. Tektronix power analyzer retains high accuracy even at very low power factor, which is very important for product characterization and development. 8.5 Crest Factor Tektronix power analyzers can measure a high Crest Factor (~10). This is critical for characterization of switched mode power supplies, which usually have high peak current draw. It has already been shown that for a sinusoidal waveform: The relationship between peak and RMS is known as the crest factor and is defined as: Thus, for a sinusoid: Many items of modern equipment connected to the AC supply take non-sinusoidal current waveforms. These include power supplies, lamp dimmers, and even fluorescent lamps. Power supplies can often exhibit a current crest factor of around 4 and up to 10. 8.6 Harmonic Distortion If a load introduces distortion of the current waveform, it is useful, in addition to knowing the crest factor, to quantify the level of distortion of the wave shape. Observation on an oscilloscope will indicate distortion but not the level of distortion. It can be shown by Fourier analysis that a nonsinusoidal current waveform consists of a fundamental component at the supply frequency plus a series of harmonics (i.e.) components at frequencies that are integral multiples of the supply frequency). For example, a SMPS, a lamp dimmer or even a speed-controlled washing machine motor can contain harmonics of even greater significance as shown in Figure 25. Figure 25. An example of Harmonics Barchart. The only useful current is the fundamental component of current, as it is only this that can generate useful power. The additional harmonic current not only flows within the power supply itself, but in all of the distribution cables, transformers and switchgear associated with the power supply and will thus cause additional loss. There is an increasing awareness of the need to limit the level of harmonics that equipment can produce. Controls exist in many territories to provide mandatory limits on the level of harmonic current permitted for certain types of load. Such regulatory controls are becoming more widespread with the use of internationally recognized standards such as EN61000-3-2. Thus, there is a need for an increased awareness amongst equipment designers as to whether their products generate harmonics and at what level. It is only the fundamental which generates power, harmonics in general do not. 8.7 Standby Power Standby power is the power drawn by a power supply when its load is not performing its full function. This may be the power consumed just by the clock on a microwave oven or the power drawn by a laptop charger when the battery is fully charged. To make the measurement requires not only very low measurement ranges, but special techniques to overcome the problems of power supplies operating in burst mode. Tektronix power analyzers have both a quick 1-button standby mode for the designer and, together with PWRVIEW pc software will perform full compliance ENERGY STAR and IEC62310 Ed.2 standby power measurements. Please see our detailed application note on this subject. 9. Conclusion Power and power related measurements on power supplies require sophisticated and accurate instrumentation to ensure that the power supply performs to its specification. The Tektronix PA1000 power analyzer incorporates a wide range of advanced features that provide class-leading accuracy from mW to kW. With the PA1000, Power supply designers faced with the challenges of ever increasing efficiency and lower standby power requirements can always be confident that their designs meet specifications.
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Houghton corrector: telescopeѲptics.net ▪ ▪ ▪ ▪ ▪▪▪▪ ▪ ▪ ▪ ▪ ▪ ▪ ▪ ▪ ▪ CONTENTS ◄ 10.2.4.3. Houghton telescopes ▐ 10.2.4.6. Plano-symmetrical HCT ► 10.2.4.5. Houghton corrector: secondary spectrum reduction As the ray plots clearly show, chromatism of the aplanatic (symmetrical) single-glass Houghton corrector becomes unacceptably large as the mirror relative aperture approaches and exceeds ~ƒ/3. While both, secondary spectrum and spherochromatism contribute to it, it is the former that dominates. The cause of the Houghton secondary spectrum is that its lens elements are not ideal thin lenses with near-zero thickness and separation. Consequently, the thin lens focal length formula given by (1/ƒ)=(1/ƒ1)+(1/ƒ2) doesn't apply. Instead, the effective focal length is given by the thick lens focal length formula for a pair of lenses, with the individual focal lengths of the front and rear lens, as the actual thick-lens focal lengths, given by: respectively, with the quantity Π being the separation between the second principal plane of the front lens and the first principal plane of the rear lens (FIG. 202) given by: with ƒ1=(n-1)[(1/R1)-1/R2] and ƒ2=(n-1)[(1/R3)-1/R4] being the front and rear lens focal length, respectively, t1 and t2 the respective lens center thicknesses, s the lens separation and R1, R2, R3 and R4 the respective lens surface radii. With 1/ƒ1 and 1/ƒ2 being very close numerically, it is the value of Π that mainly determines corrector's power (Eq. 150), and hence its secondary spectrum. FIGURE 202: Houghton corrector principal plane separation Π, between the second principal plane of the front lens (P1') and the first principal plane of the second lens (P2), is the main determinant of the effective non-zero focal length of the doublet. According to Eq. 151, the principal plane separation increases with lens thickness and lens spacing (s), reducing further corrector (positive) focal length, nearly in proportion to ~ƒ1ƒ2/Π. This in turn results in the increase of corrector's secondary spectrum. The principal plane separation - and secondary spectrum - for given aperture also increase with more strongly curved lens surfaces, needed to correct spherical aberration and coma of faster mirrors. When the two lenses are in contact, the separation s is likely to be, roughly, somewhat less than a half of the value of Π. This means that any relative increase in lens separation will result in less than half as much of a relative increase of secondary spectrum. With the typical values of ƒ1/R1~1.3 and ƒ2/R4~0.5, change in the thickness of the front lens results in secondary spectrum larger by a factor of ~2.5 than with identical thickness increase of the rear lens. Also, any relative change in the front lens thickness will result in ~1.3 times greater secondary spectrum than identical increase in lens separation, while thickness increase of the rear lens will result in the increase of secondary spectrum smaller by a factor of 0.5 than identical increase in lens separation. As can be seen from Eq. 151, the size of Π it is not affected by a change in index, or radii, but it does change in proportion to lens thickness. Since the lenses can only be so thin, there is a practical minimum level of secondary spectrum that can not be further lessened with a single-glass Houghton corrector. It could only be cancelled by changing the sign of R4, but it would require the rear element to be a negative out-curving meniscus, which would make corrector with acceptable coma impossible. In the symmetrical corrector configuration, the positive element has slightly less power than the negative element as a single lens, but when combined the two have a weak positive power. Neglecting the smaller (negative) value of (1/ƒ1)-1/ƒ2 in Eq. 150, the corrector focal length is approximated by ƒc~ƒ1ƒ2/Π (at the end of this section is explained in more details why taking somewhat smaller value is likely to give better end result), with the variation in the focal length for non-optimized wavelengths approximated by where n and n' are the refractive index of the optimized and non-optimized wavelength, respectively. It places the secondary spectrum of the Houghton corrector - as the axial separation of its paraxial foci - roughly at ~ƒc/200, or some 10 times greater than in a doublet achromat. It is only due to the weak corrector's power that it induces relatively small amount of chromatism. Nominally, the secondary spectrum is negative for shorter (than the optimized) wavelengths, and positive for longer wavelengths - a consequence of the positive effective power of the corrector. It simply means that the former focus shorter, and the latter longer than the optimized wavelength. Should the secondary spectrum be the only chromatic aberration present, the transverse chromatic blur diameter in units of the green e-line Airy disc diameter would be approximated by B~745(Δƒ)/Fc2, with Fc being the corrector's focal ratio, approximated by Fc~ƒc/D~ƒ1ƒ2/ΠD. In reality, secondary spectrum is always combined with a certain amount of sphero-chromatism, in which case the actual secondary spectrum is measured by the separation between best aberrated foci for different wavelengths, and the blur size results from the combined size of defocus (secondary spectrum) and spherical aberration (sphero-chromatism) for the wavelength. Houghton corrector secondary spectrum can be minimized by either slightly weakening chromatic power of the front lens, or by slightly strengthening the rear lens. The two options for reducing the secondary spectrum of the Houghton corrector are: (1) abandoning symmetrical radii, and (2) using two different glass types for the lenses. Another option is to use plano-symmetrical corrector type, which has generally less of the residual power, hence smaller secondary spectrum. Abandoning equal radii design is a practical disadvantage fabrication-wise (also requiring more complex calculations), while plano-lens design puts constraints on coma-correction. Thus, the most effective way of reducing the Houghton secondary spectrum is by using two different glass types. Starting out with a typical single-glass aplanatic Houghton corrector, minimizing the secondary spectrum requires small change in either power, or dispersion (or both) of one of the two lens elements. Lens power changes with (n-1), and its dispersion with 1/V, n being the glass refractive index and V its nominal dispersion. The measure of needed chromatic power change can be obtained from the general rule of achromatism for near-contact or contact doublet, requiring ƒ1/ƒ2=V2/V1. Since the aplanatic Houghton doublet acts as a thin lens pair in which the rear lens focal length is somewhat longer than that of the front lens, resulting in a weak positive power of the doublet, the appropriate dispersion V2 for minimizing the secondary spectrum should be different in approximately the same proportion. Assuming thin-lens doublet, the rear lens effective focal length ƒ2e is obtained from 1/ƒc=(1/ƒ1)+1/ƒ2e, as 1/ƒ2e=1/ƒc-(1/ƒ1). The appropriate rear lens dispersion V2 can be obtained from ƒ1/ƒ2e=V2/V 0, with V0 being the dispersion of the basic single-glass corrector. However, glass of different dispersion will almost invariably also have different refractive index. Since it is also a factor affecting chromatic power (secondary spectrum), it needs to be taken into account. With it, the achromatic relation takes the form ƒ1/ƒ2e=(n0-1)V2/(n2-1)V0, with needed properties of the achromatizing glass obtained from: with n0, V0 and n2, V2 being the refractive index (optimized wavelength) and Abbe number for the starting single-glass aplanatic corrector, and the achromatized negative element, respectively. Helpful indicator of the effect of new glass type on secondary spectrum is given by 1/ç=(n2-1)V0/(n0-1)V2. This parameter can be called corrector's relative chromatic power. If it is the front lens to be achromatized, ç will be slightly smaller, and for the rear (negative) lens slightly greater than 1 (rough average with n~1.5 and mirrors in the ƒ/2.5-ƒ/3 range is ±1%). It reflects needed change in the relative chromatic power of one of the lens elements for near-optimal balancing of the combined secondary spectrum of the corrector. In general, the second glass alternative for minimizing the secondary spectrum will be of very similar index and dispersion values to those of the single-glass corrector to be achromatized (as the nominal dispersion and index change in opposite directions with significant index changes, it becomes difficult to impossible to find a glass with relative chromatism ç close to 1). Once suitable glass is found, it will make possible significant reduction of the secondary spectrum. An integral part of the optimization is tailoring out best combination of the secondary spectrum and spherical aberration, with lens thickness and separation also being possible factors (FIG. 203). FIGURE 203: Effect of achromatizing on the chromatism of a HCT with symmetrical aplanatic Houghton corrector (compare with Fig. 132a and 133a). By replacing the BK7 negative lens element glass with BK8, the wavefront error is reduced from 1.3 and 0.38 to 0.45 and 0.048 wave RMS for the h and r spectral lines, respectively (system a in the Appendix). By increasing lens separation, the chromatism is further reduced and optimized (balanced) to 0.14 and 0.12 wave RMS for the h and r lines, respectively, as shown on the plot (system b in the Appendix). The h-line correction is now nearly twice better than in non-achromatized plano-lens version, or at the level of a 4" ƒ/200 achromat. While still approximately double the chromatism of a comparable SCT, it is vastly improved over the non-achromatized version (note that further reduction could be possible). An ƒ/2.5 relative aperture for this aperture size is probably near-limit for all-spherical Houghton system. At this primary focal ratio, higher-order spherical, responsible for most of the scattered rays, cannot be corrected significantly better than 1/30 wave RMS without adding aspheric surface term. SPEC'S Achromatizing is still worthwhile at ~ƒ/3 mirror focal ratios. By switching a rear lens from BK7 to PK3, and increasing lens spacing from 1.4mm to 3mm (to compensate for the induced spherical aberration) in the ƒ/3/10 symmetrical aplanatic Houghton (Fig. 133c), corrector's chromatism is reduced from 0.37 and 0.1 to 0.085 and 0.043 wave RMS wavefront error for the violet h-line and red r-line, respectively. That compares very favorably to the reduction of chromatism by compromising it with some coma (0.28 and 0.075 wave RMS for the h- and r-lines, respectively). Note that the corrector type with all four radii different allows for still better color correction in a single-glass arrangement. An alternative to increasing lens separation of the achromatized corrector is slightly (typically less than 1%) relaxing second radius. It is more likely to bring best foci of different wavelengths closer together, but the choice is best made using ray-trace. As mentioned, secondary spectrum is main, but not a lone contributor to the chromatism of the Houghton corrector. The other is sphero-chromatism. As a consequence, the secondary spectrum - and chromatism in general - is minimized by bringing together best foci for different wavelengths, not the paraxial foci. This would certainly require more involved optimization than the one outlined above. In general, however, considering typical symmetrical aplanatic Houghton LA properties, even this crude form of optimization alone - in particular with the change in relative chromatic power purposely reduced to ~2/3 of the power differential indicated by the ç value - should result in a significant reduction of corrector's chromatism. The achromatizing glass often will, to some extent, also change the spherochromatism of the basic (single-glass) corrector. It can be for better, or for worse, but the effect is, in general, secondary to that of the change in size of secondary spectrum. ◄ 10.2.4.3. Houghton telescopes ▐ 10.2.4.6. Plano-symmetrical HCT ► Home | Comments
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ACO Seminar The ACO Seminar (2014–2015) Jan. 15, 3:30pm, Wean 8220 Alexander Barvinok, University of Michigan Computing partition functions in hard problems of combinatorial optimization Consider an NP-complete problem, such as finding whether a given graph vertices contains a Hamiltonian cycle. First, we attempt to solve an even more general problem: given an non-negative matrix A, interpreted as the matrix of weights on the edges of the complete graph , we consider the sum (partition function) over all Hamiltonian cycles in of the products of weights on the edges of the cycle. Thus if the adjacency matrix of , we obtain the number of Hamiltonian cycles in . Next, we modify by replacing all zeros with some positive epsilons. It turns out that the partition function becomes efficiently computable, which allows us to distinguish graphs that are far from having a Hamiltonian cycles from graphs that have many Hamiltonian cycles (even when "many" still means that the probability to hit a Hamiltonian cycle at random is exponentially small). In a joint work with P. Soberon, we attempt to establish the most general framework when this approach appears to work: it concerns computing the partition function of graph homomorphisms with multiplicities, which allows us to handle Hamiltonian cycles, independent sets, dense subgraphs and colorings.
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10000 Mortgage Over 15 Years 10000 Mortgage Over 15 Years For those looking to secure a Texas land mortgage, our land loan calculator Input your loan term (total years on the loan). Determine your payment. Most common terms are 15 years and 30 years. Mortgage amount: Total balance Annual interest rate for this mortgage without purchasing any discount points. For example, if you're interested in paying off your mortgage off in 15 years as opposed to 30, you generally need a monthly payment that is X your typical. Loan option. Select your mortgage term length *. Fixed 30 Years, Fixed 20 Years, Fixed 15 Years Total payments over 30 years. $XXXX. Principal. Interest. Mortgage. Get the edge over other buyers, and a $10, seller guarantee. I How Much House Can I Afford? Refinance Break Even Calculator · 30 to 15 Year. # of Payments is the number of monthly payments you will make to pay off the loan. For example, if the approximate term of the loan is 4 years or 48 months, you. 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SciPy Programming Succinctly - Free Computer, Programming, Mathematics, Technical Books, Lecture Notes and Tutorials SciPy Programming Succinctly • Title SciPy Programming Succinctly • Author(s) James McCaffrey • Publisher: Syncfusion (2016) • Hardcover/Paperback N/A • eBook HTML, PDF (120 pages), ePub, adn MOBI (Kindle) • Language: English • ISBN-10: N/A • ISBN-13: N/A • Share This: Book Description This book offers readers a quick, thorough grounding in knowledge of the Python open source extension SciPy. The SciPy library, accompanied by its interdependent NumPy, offers Python programmers advanced functions that work with arrays and matrices. Each section presents a complete demo program for programmers to experiment with, carefully chosen examples to best illustrate each function, and resources for further learning. Use this e-book to install and edit SciPy, and use arrays, matrices, and combinatorics in Python programming. About the Authors Reviews, Ratings, and Recommendations: Related Book Categories: Read and Download Links:Similar Books: Book Categories Other Categories Resources and Links
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Eilenberg-Steenrod Axioms -- from Wolfram MathWorld A family of functors category of pairs of topological spaces and continuous maps, to the category of Abelian groups and group homomorphisms satisfies the Eilenberg-Steenrod axioms if the following conditions hold. 1. long exact sequence of a pair axiom. For every pair where the map inclusion map inclusion map map boundary map. 2. homotopy axiom. If induced maps 3. excision axiom. If space with subspaces set closure of inclusion map 4. dimension axiom. Let groups. The coefficients of the homology theory These are the axioms for a generalized homology theory. For a cohomology theory, instead of requiring that functor, it is required to be a co-functor (meaning the induced map points in the opposite direction). With that modification, the axioms are essentially the same (except that all the induced maps point backwards).
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Understanding the Area of a Cube Understanding the Area of a Cube Introduction to the Area of a Cube Cubes are geometric marvels that we encounter in everyday life, from dice in our game nights to shipping boxes. But beyond their boxy charm lies an interesting mathematical concept: their surface area. Calculating the area of a cube is a fundamental concept in geometry that provides valuable insights for various real world applications. Let’s dive into it! Dissecting the Formula The formula to find the area of a cube is simple yet powerful: A = 6s². Here: • A represents the total surface area of the cube, expressed in square units like square meters (m²) or square feet (ft²). • s is the length of one side of the cube, expressed in linear units like meters (m) or feet (ft). In essence, the surface area (A) is equal to six times the square of the side length (s). Real life Example: Packaging Design Imagine you are designing a gift box for a new product launch. You’ve settled on a chic cube shaped box with each side measuring 0.5 meters. What’s the total surface area? Plugging into the formula, we have: A = 6 * (0.5)² = 6 * 0.25 = 1.5 m² Thus, you’ll need 1.5 square meters of material to cover the entire surface of the cube. Practical Application: Construction Engineers and architects regularly use this formula in designing structures. For instance, if a company plans to construct cube shaped storage units, knowing the surface area helps in estimating material costs. Data Validation and Practical Limitations It’s important to ensure that the side length (s) is a positive number. Negative or zero values are not physically meaningful for length and should return an error message. Calculating the area of a cube is a straightforward yet invaluable skill in geometry. From packaging design to construction, this formula A = 6s² helps you quantify the surface area required for various practical applications. Understanding this basic formula opens the door to numerous real world applications, making it an essential tool in both education and industry. Q: Can the side length (s) of a cube be in different units? A: Yes, the side length can be in any linear unit like meters, feet, inches, etc. Just ensure consistency when calculating the area. Q: What if the side length is zero or negative? A: The side length should be a positive number. Zero or negative values don’t make sense and should return an error message. Example Calculations 1. s = 1 m Surface Area: A = 6 * 1² = 6 m² 2. s = 2 ft Surface Area: A = 6 * 2² = 24 ft² 3. s = 3 cm Surface Area: A = 6 * 3² = 54 cm² Tags: Geometry, Mathematics, Cube
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Andre Weil: Founding Member of the Mathematical Bourbaki Group ANDRE WEIL: Founding Member of the Mathematical Bourbaki Group André Weil was a very influential French mathematician around the middle of the 20th Century. Born into a prosperous Jewish family in Paris, he was brother to the well-known philosopher and writer Simone Weil, and both were child prodigies. He was passionately addicted to mathematics by the age of ten, but he also loved to travel and study languages (by the age of sixteen he had read the “Bhagavad Gita” in the original Sanskrit). He studied (and later taught) in Paris, Rome, Göttingen and elsewhere, as well as at the Aligarh Muslim University in Uttar Pradesh, India, were he further explored what would become a life-long interest in Hinduism and Sanskrit literature. Even as a young man, Weil made substantial contributions in many areas of mathematics, and was particularly animated by the idea of discovering profound connections between algebraic geometry and number theory. His fascination with Diophantine equations led to his first substantial piece of mathematical research on the theory of algebraic curves. During the 1930s, he introduced the adele ring, a topological ring in algebraic number theory and topological algebra, which is built on the field of rational numbers. The early leader of the Bourbaki group Weil was an early leader of the Bourbaki group who published many influential textbooks on modern mathematics It was also at this time that he became a founding member, and the de facto early leader, of the so-called Bourbaki group of French mathematicians. This influential group published many textbooks on advanced 20th Century mathematics under the assumed name of Nicolas Bourbaki, in an attempt to give a unified description of all mathematics founded on set theory. Bourbaki has the distinction of having been refused membership of the American Mathematical Society for being non-existent (although he was a member of the Mathematical Society of France!) When the Second World War broke out, Weil, a committed conscientious objector, fled to Finland, where he was mistakenly arrested as a possible spy. Having made his way back to France, he was again arrested and imprisoned as for refusing to report for military service. In his trial, he cited the Bhagavad Gita to justify his stand, arguing that his true dharma was the pursuit of mathematics, not assisting in the war effort, however just the cause. Given the choice of five more years in prison or joining a French combat unit, though, he chose the latter, an especially lucky decision given that the prison was blown up shortly afterwards. But it was in 1940, in a prison near Rouen, that Weil did the work that really made his reputation (although his full proofs had to wait until 1948, and even more rigorous proofs were supplied by Pierre Deligne in 1973). Building on the prescient work of his countryman Évariste Galois in the previous century, Weil picked up the idea of using geometry to analyze equations, and developed algebraic geometry, a whole new language for understanding solutions to equations. Weil Conjectures An illustration of the “cycle évanescent” or “vanishing cycle” described in Deligne’s proof of the Weil conjectures The Weil conjectures on local zeta-functions effectively proved the Riemann hypothesis for curves over finite fields, by counting the number of points on algebraic varieties over finite fields. In the process, he introduced for the first time the notion of an abstract algebraic variety and thereby laid the foundations for abstract algebraic geometry and the modern theory of abelian varieties, as well as the theory of modular forms, automorphic functions and automorphic representations. His work on algebraic curves has influenced a wide variety of areas, including some outside of mathematics, such as elementary particle physics and string theory. In 1941, Weil and his wife took the opportunity to sail for the United States, where they spent the rest of the War and the rest of their lives. In the late 1950s, Weil formulated another important conjecture, this time on Tamagawa numbers, which remained resistant to proof until 1989. He was instrumental in the formulation of the so-called Shimura-Taniyama-Weil conjecture on elliptic curves which was used by Andrew Wiles as a link in the proof of Fermat’s Last Theorem. He also developed the Weil representation, an infinite-dimensional linear representation of theta functions which gave a contemporary framework for understanding the classical theory of quadratic forms. Over his lifetime, Weil received many honorary memberships, including the London Mathematical Society, the Royal Society of London, the French Academy of Sciences and the American National Academy of Sciences. He remained active as professor emeritus at the Institute for Advanced Studies at Princeton until a few years before his death. << Back to Turing Forward to Cohen >>
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Varney, Annisa Varney, Anthea Varney, Bethel Varney, Charlesetta Varney, Davetta Varney, Deonne Varney, Emmie Varney, Fanta Varney, Felicitas Varney, Flavia Varney, Glenn Varney, Gwendolynn Varney, Hollye Varney, Jalene Varney, Jamel Varney, Jenita Varney, Jonda Varney, Kameka Varney, Lagina Varney, Lateasha Varney, Louanne Varney, Ludivina Varney, Ramonita Varney, Reshonda Varney, Romina Varney, Samone Varney, Shauntae Varney, Shelle Varney, Tasheka Varney, Tashima Varney, Tekesha Varney, Trini Varney, Tyese Varney, Alesa Varney, Brigett Varney, Camella Varney, Charee Varney, Coby Varney, Constantina Varney, Contrina Varney, Danise Varney, Delanie Varney, Denia Varney, Detrice Varney, Donyelle Varney, Estee Varney, Freya Varney, Gabriele Varney, Heidie Varney, Idella Varney, Iraida Varney, Jonica Varney, Kaley Varney, Kirsty Varney, Laurette Varney, Nataki Varney, Patrisia Varney, Ronnetta Varney, Signe Varney, Tarita Varney, Trenia Varney, Annessa Varney, Celene Varney, Cerissa Varney, Chinita Varney, Danah Varney, Dannelle Varney, Dawnielle Varney, Doretta Varney, Dwanna Varney, Jinger Varney, Karilyn Varney, Loida Varney, Miko Varney, Mirtha Varney, Myia Varney, Perri Varney, Ragina Varney, Rea Varney, Saralyn Varney, Saskia Varney, Sharry Varney, Tiare Varney, Tonica Varney, Waleska Varney, Xochilt Varney, Zipporah Varney, Delynn Varney, Detria Varney, Earl Varney, Enedelia Varney, Jeanett Varney, Jonita Varney, Kenita Varney, Laural Varney, Leonard Varney, Ligia Varney, Linnie Varney, Lowanda Varney, Ltanya Varney, Monquie Varney, Paloma Varney, Rafael Varney, Steve Varney, Torina Varney, Yolunda Varney, Cariann Varney, Carmita Varney, Cerise Varney, Dalana Varney, Dawnya Varney, Deshannon Varney, Deshonda Varney, Desirea Varney, Eartha Varney, Emelia Varney, Francoise Varney, Gerrie Varney, Gitty Varney, Jared Varney, Kianna Varney, Lamanda Varney, Laraine Varney, Latonda Varney, Latrease Varney, Leaann Varney, Liv Varney, Maha Varney, Meliza Varney, 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Melessia Varney, Nicolasa Varney, Quantina Varney, Remy Varney, Robina Varney, Roya Varney, Saba Varney, Sacheen Varney, Shalee Varney, Sharlet Varney, Shereese Varney, Sherronda Varney, Shewanda Varney, Skyla Varney, Tallie Varney, Tanganyika Varney, Taressa Varney, Tarshia Varney, Timi Varney, Zola Varney, Ainsley Varney, Alivia Varney, Ardith Varney, Artina Varney, Ashely Varney, Athanasia Varney, Chantale Varney, Chimene Varney, Corinn Varney, Daffney Varney, Danene Varney, Daphnie Varney, Deeanne Varney, Erynn Varney, Glennis Varney, Kalee Varney, Kandyce Varney, Kathren Varney, Kelcey Varney, Laticha Varney, Lenee Varney, Luvenia Varney, Machele Varney, Maris Varney, Marisella Varney, Nakeia Varney, Nickey Varney, Nikkita Varney, Noella Varney, Peri Varney, Raushanah Varney, Renisha Varney, Shahidah Varney, Shakisha Varney, Shevonne Varney, Tashema Varney, Tondalaya Varney, Trinita Varney, Akiko Varney, Alfredia Varney, Bathsheba Varney, Billye Varney, Carter Varney, Clementina Varney, Danyele Varney, Darnisha Varney, Darryl Varney, Eda Varney, Emy Varney, Faustina Varney, Hayden Varney, Hortensia Varney, Ivon Varney, Javonna Varney, Keita Varney, Lakeisa Varney, Lan Varney, Leota Varney, Letasha Varney, Maggi Varney, Monita Varney, Nadirah Varney, Nikka Varney, Noemy Varney, Roxy Varney, Sabrinia Varney, Sanda Varney, Sharan Varney, Sharnette Varney, Swati Varney, Tiny Varney, Tonjia Varney, Trella Varney, Venecia Varney, Zana Varney, Zorana Varney, Adrienna Varney, Ama Varney, Angelisa Varney, Beckey Varney, Bobbye Varney, Cammi Varney, Chequita Varney, Clinton Varney, Dama Varney, Danee Varney, Dari Varney, Deandre Varney, Deetta Varney, Devonda Varney, Edelmira Varney, Ieasha Varney, Joylynn Varney, Krisinda Varney, Kristol Varney, Kwana Varney, Lacresia Varney, Ladawna Varney, Lavanda Varney, Lavena Varney, Leaha Varney, Massiel Varney, Michiko Varney, Naketa Varney, Persephone Varney, Raynell Varney, Ronell Varney, Shalynn Varney, Shimeka Varney, 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Varney, Cartina Varney, Cathlene Varney, Chaquita Varney, Clifford Varney, Danitra Varney, Dayana Varney, Delania Varney, Dennie Varney, Edina Varney, Electra Varney, Elouise Varney, Ginelle Varney, Heath Varney, Kassia Varney, Katrenia Varney, Keidra Varney, Laconda Varney, Loura Varney, Marchell Varney, Mayela Varney, Merle Varney, Merlinda Varney, Mychelle Varney, Nannie Varney, Nikko Varney, Nori Varney, Prescilla Varney, Rika Varney, Ronika Varney, Sahara Varney, Sandrea Varney, Senora Varney, Shawntee Varney, Sheryll Varney, Shina Varney, Terre Varney, Tinna Varney, Tonnette Varney, Torre Varney, Tya Varney, Vianey Varney, Zora Varney, Adonna Varney, Alberto Varney, Almee Varney, Ashanta Varney, Bridgid Varney, Cindia Varney, Cricket Varney, Darah Varney, Dineen Varney, Elan Varney, Iman Varney, Jacquelynne Varney, Jennilyn Varney, Jessicca Varney, Julio Varney, Kenosha Varney, Kimyata Varney, Lanessa Varney, Leesha Varney, Letoya Varney, Nkenge Varney, Priti Varney, Raizel Varney, Rosaisela Varney, Sameerah Varney, Sera Varney, Shallon Varney, Sharmane Varney, Shauntay Varney, Sherhonda Varney, Shonnie Varney, Sueanne Varney, Taneesha Varney, Telina Varney, Tenecia Varney, Tyanne Varney, Alexsandra Varney, Angeles Varney, Charlean Varney, Danel Varney, Davi Varney, Daysha Varney, Demaris Varney, Dorenda Varney, Dorita Varney, Dorlisa Varney, Genee Varney, Jamella Varney, Kathyjo Varney, Kaya Varney, Latondra Varney, Leshia Varney, Mahala Varney, Marija Varney, Maudie Varney, Megann Varney, Phuong Varney, Reesa Varney, Ronya Varney, Selinda Varney, Shama Varney, Shirell Varney, Shonita Varney, Taiwana Varney, Takenya Varney, Talonda Varney, Tassie Varney, Telena Varney, Theodore Varney, Violetta Varney, Willetta Varney, Adora Varney, Alvin Varney, Chauncey Varney, Chrisy Varney, Clancy Varney, Claressa Varney, Corrinna Varney, Darra Varney, Deshawna Varney, Donise Varney, Elona Varney, Gara Varney, Georgene Varney, Gila Varney, Jania Varney, Jocelynn Varney, Joei Varney, Kassy Varney, Liset Varney, Maleka Varney, Mashell Varney, Melida Varney, Michalene Varney, Moya Varney, Nyisha Varney, Rayanne Varney, Renia Varney, Salima Varney, Samanthia Varney, Sammy Varney, Santosha Varney, Shaneeka Varney, Shawntelle Varney, Shwanda Varney, Sita Varney, Tarin Varney, Tawn Varney, Tosca Varney, Valisha Varney, Vangie Varney, Vernon Varney, Zella Varney, Amarilis Varney, Amelie Varney, Aubri Varney, Bobbette Varney, Byron Varney, Chance Varney, Deshon Varney, Eraina Varney, Jaycee Varney, Joyell Varney, Keirsten Varney, Kendi Varney, Kendrea Varney, Keyanna Varney, Khristie Varney, Kimberlyann Varney, Kimyetta Varney, Kinberly Varney, Lakeita Varney, Lakina Varney, Leida Varney, Lenae Varney, Mariette Varney, Marymargaret Varney, Meera Varney, Micky Varney, Naja Varney, Oneka Varney, Pedro Varney, Rennie Varney, Sharlee Varney, Sharonne Varney, Shonya Varney, Solange Varney, Tashawn Varney, Amorette Varney, Andrienne Varney, Ari Varney, 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Varney, Takeysha Varney, Tamelia Varney, Tanyika Varney, Temperance Varney, Theodosia Varney, Thyra Varney, Tyria Varney, Urania Varney, Yvett Varney, Angelynn Varney, Anitria Varney, Anny Varney, Arlisha Varney, Aya Varney, Brielle Varney, Careen Varney, Correne Varney, Dagmar Varney, Dalynn Varney, Dannetta Varney, Delight Varney, Denetta Varney, Desaree Varney, Donyel Varney, Eisha Varney, Eleanora Varney, Felix Varney, Fransisca Varney, Jannel Varney, Jenia Varney, Jettie Varney, Jolina Varney, Jowanna Varney, Junko Varney, Kaaren Varney, Kathey Varney, Keryn Varney, Ketina Varney, Korrine Varney, Kristalyn Varney, Laini Varney, Lamona Varney, Lanay Varney, Leza Varney, Maddalena Varney, Mellody Varney, Merica Varney, Merita Varney, Nafeesa Varney, Nicholl Varney, Oneika Varney, Pamelyn Varney, Romonda Varney, Ronique Varney, Rosslyn Varney, Rosy Varney, Sachiko Varney, Sarajane Varney, Satara Varney, Shamra Varney, Tamberly Varney, Tawona Varney, Tiasha Varney, Tine Varney, 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Varney, Vanya Varney, Verdell Varney, Vonna Varney, Zsazsa Varney, Adonia Varney, Alease Varney, Aleena Varney, Alishea Varney, Anetta Varney, Ania Varney, Anquinette Varney, Antonieta Varney, Archana Varney, Arely Varney, Artavia Varney, Arturo Varney, Batya Varney, Bernardine Varney, Betti Varney, Breck Varney, Camara Varney, Cesar Varney, Chalise Varney, Charly Varney, Cherrise Varney, Chikita Varney, Cristiana Varney, Dayatra Varney, Ericha Varney, Ermalinda Varney, Eulonda Varney, Fayth Varney, Halina Varney, Hunter Varney, Iveliz Varney, Ivie Varney, Jahna Varney, Jancy Varney, Janisha Varney, Jeaninne Varney, Jeannemarie Varney, Jessamine Varney, Keeya Varney, Kelisha Varney, Kenyette Varney, Kyley Varney, Kyrie Varney, Ladora Varney, Laguana Varney, Larri Varney, Leyda Varney, Louie Varney, Marykathryn Varney, Mekesha Varney, Meredeth Varney, Merridith Varney, Mittie Varney, Nakeitha Varney, Narda Varney, Nature Varney, Neeta Varney, Nitasha Varney, Nuvia Varney, Paticia Varney, Promise Varney, Rabia Varney, Raquelle Varney, Richetta Varney, Shawneequa Varney, Shervon Varney, Siomara Varney, Starlyn Varney, Tahirih Varney, Taia Varney, Talita Varney, Tamina Varney, Tamya Varney, Tanicka Varney, Tanyanika Varney, Tawan Varney, Tawanya Varney, Tiarra Varney, Timisha Varney, Tinita Varney, Tonnia Varney, Toronda Varney, Toshika Varney, Triva Varney, Trula Varney, Tykesha Varney, Vaness Varney, Vennessa Varney, Wakeelah Varney, Wilhemina Varney, Yoshiko Varney, Afia Varney, Afiya Varney, Allysa Varney, Alsha Varney, Amberlyn Varney, Arah Varney, Audri Varney, Bernardette Varney, Bernedette Varney, Binta Varney, Breeann Varney, Britni Varney, Carnetta Varney, Cesilia Varney, Chelli Varney, Dawnna Varney, Desirie Varney, Fleur Varney, Gianina Varney, Hanh Varney, Ivelis Varney, Jalena Varney, Jenafer Varney, Johann Varney, Joyanne Varney, Kamia Varney, Katerine Varney, Kecha Varney, Keir Varney, Kennita Varney, Kimra Varney, Kristiann Varney, Levita Varney, Lluvia Varney, Lorian Varney, Lyndsy Varney, Mardell Varney, Maurissa Varney, Mischell Varney, Mistey Varney, Nataya Varney, Nickola Varney, Obdulia Varney, Preston Varney, Raynetta Varney, Rondell Varney, Shakara Varney, Shalan Varney, Shanya Varney, Shawnika Varney, Shenetha Varney, Shilah Varney, Soo Varney, Tametha Varney, Tanicia Varney, Tanny Varney, Tennie Varney, Tiffannie Varney, Torra Varney, Trinika Varney, Val Varney, Abeer Varney, Adah Varney, Aleece Varney, Andreya Varney, Annaliza Varney, Arrie Varney, Ayasha Varney, Berit Varney, Birgit Varney, Bre Varney, Brigetta Varney, Cassundra Varney, Christyl Varney, Cristyn Varney, Danean Varney, Daralyn Varney, Davonne Varney, Deisha Varney, Denina Varney, Dietrich Varney, Domini Varney, Erinne Varney, Evy Varney, Forrest Varney, Genetta Varney, Gioia Varney, Guenevere Varney, Gwenetta Varney, Gwenette Varney, Hesper Varney, Ignacia Varney, Irmalinda Varney, Jacquita Varney, Jeanita Varney, Jodiann Varney, Kalin Varney, Kenny Varney, Keon Varney, Kimberlynn Varney, Kimela Varney, Kimm Varney, Kiwanna Varney, Lacreshia Varney, Lakashia Varney, Latarshia Varney, Latrish Varney, Lauria Varney, Likisha Varney, Lindee Varney, Little Varney, Loreta Varney, Lutisha Varney, Mahasin Varney, Maki Varney, Maxi Varney, Meesha Varney, Melitta Varney, Mery Varney, Mysty Varney, Quanisha Varney, Ranetta Varney, Rosia Varney, Sadonna Varney, Sangita Varney, Sanora Varney, Shanina Varney, Sharonna Varney, Shavona Varney, Sherrica Varney, Storm Varney, Tabrina Varney, Tacha Varney, Taleshia Varney, Tamee Varney, Tanetta Varney, Thanh Varney, Twanya Varney, Ulrica Varney, Yaffa Varney, Aarin Varney, Afi Varney, Aliscia Varney, Aneta Varney, Anetria Varney, Antoinett Varney, Athenia Varney, Atisha Varney, Barb Varney, Brit Varney, Caley Varney, Camala Varney, Carena Varney, Carinne Varney, Carmine Varney, Charlott Varney, Chessa Varney, Chundra Varney, Cotrina Varney, Drenda Varney, Elanor Varney, Eleftheria Varney, Erricka Varney, Feleshia Varney, Geanna Varney, Harley Varney, Honora Varney, Jakia Varney, Jannah Varney, Jennene Varney, Jennett Varney, Jobeth Varney, Jordanna Varney, Juanetta Varney, Junie Varney, Kaira Varney, Kashana Varney, Katee Varney, Kateena Varney, Kelleigh Varney, Kenyana Varney, Kerrilyn Varney, Kessa Varney, Kirk Varney, Lakitha Varney, Laneshia Varney, Lanissa Varney, Lashunta Varney, Leonia Varney, Lester Varney, Letina Varney, Llesenia Varney, Loranda Varney, Lorilynn Varney, Lotus Varney, Luke Varney, Marisal Varney, Marvis Varney, Maurie Varney, Minnette Varney, Mishawn Varney, Natsha Varney, Ngoc Varney, Nioka Varney, Particia Varney, Phillina Varney, Prisilla Varney, Raney Varney, Raphaela Varney, Rashaunda Varney, Raynelle Varney, Ross Varney, Rufina Varney, Sarahann Varney, Shaconda Varney, Shanicka Varney, Shaunita Varney, Shekina Varney, Shelie Varney, Sherah Varney, Shermeka Varney, Shermika Varney, Sonnia Varney, Stephonie Varney, Taffie Varney, Takima Varney, Talea Varney, Tamaya Varney, Taneika Varney, Tanina Varney, Tanisia Varney, Taysha Varney, Tekeshia Varney, Terie Varney, Therasa Varney, Thomasa Varney, Tifini Varney, Tomie Varney, Tressy Varney, Trissa Varney, Tristy Varney, Tritia Varney, Twania Varney, Tyrene Varney, Uganda Varney, Wandy Varney, Wileen Varney, Yumeka Varney, Zarina Varney, Adreana Varney, Alandra Varney, Alenda Varney, Amera Varney, Arlicia Varney, Artemis Varney, Avani Varney, Bena Varney, Candia Varney, Cathyjo Varney, Chala Varney, Charitie Varney, Charlise Varney, Chasiti Varney, Chayla Varney, Chrisanne Varney, Christell Varney, Clarrissa Varney, Cortnee Varney, Dametria Varney, Daneka Varney, Darlisa Varney, Delorse Varney, Deshone Varney, Dorette Varney, Esperansa Varney, Evagelia Varney, Feliciana Varney, Fikisha Varney, Freeda Varney, Frida Varney, Gaila Varney, Genae Varney, Georgi Varney, Gitel Varney, Gricel Varney, Hilliary Varney, Ichelle Varney, Ishia Varney, Ivelise Varney, Jacci Varney, Jamesa Varney, Jammy Varney, Jerrilynn Varney, Jillaine Varney, Jillann Varney, Jin Varney, Joda Varney, Jodeen Varney, Kalisa Varney, Karne Varney, Katura Varney, Kayci Varney, Keegan Varney, Keenan Varney, Khia Varney, Koleen Varney, Krishana Varney, Krislyn Varney, Kurt Varney, Lafaye Varney, Lakeysa Varney, Lakresha Varney, Lalisa Varney, Lashonne Varney, Leighanna Varney, Lesslie Varney, Lin Varney, Madlyn Varney, Malky Varney, Melenda Varney, Melisssa Varney, Mikita Varney, Monicia Varney, Monik Varney, Nakima Varney, Naquita Varney, Natash Varney, Raphael Varney, Rashawna Varney, Reannon Varney, Reneta Varney, Reshanda Varney, Riana Varney, Rocky Varney, Rosamond Varney, Roshan Varney, Saleena Varney, Schuyler Varney, Shaka Varney, Shakesha Varney, Shaleta Varney, Shanekia Varney, Surina Varney, Syrita Varney, Taj Varney, Tajuanna Varney, Tasheen Varney, Tiffine Varney, Timmie Varney, Tomikia Varney, Trease Varney, Tyrhonda Varney, Wakisha Varney, Yvone Varney, Adalia Varney, Ahuva 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How hard is it to satisfy (Almost) all roommates? The classic Stable Roommates problem (the non-bipartite generalization of the well-known Stable Marriage problem) asks whether there is a stable matching for a given set of agents, i.e. a partitioning of the agents into disjoint pairs such that no two agents induce a blocking pair. Herein, each agent has a preference list denoting who it prefers to have as a partner, and two agents are blocking if they prefer to be with each other rather than with their assigned partners. Since stable matchings may not be unique, we study an NP-hard optimization variant of Stable Roommates, called Egal Stable Roommates, which seeks to find a stable matching with a minimum egalitarian cost γ, i.e. the sum of the dissatisfaction of the agents is minimum. The dissatisfaction of an agent is the number of agents that this agent prefers over its partner if it is matched; otherwise it is the length of its preference list. We also study almost stable matchings, called Min-Block-Pair Stable Roommates, which seeks to find a matching with a minimum number β of blocking pairs. Our main result is that Egal Stable Roommates parameterized by γ is fixed-parameter tractable, while Min-Block-Pair Stable Roommates parameterized by β is W[1]-hard, even if the length of each preference list is at most five. Original language English Title of host publication 45th International Colloquium on Automata, Languages, and Programming, ICALP 2018 Editors Christos Kaklamanis, Daniel Marx, Ioannis Chatzigiannakis, Donald Sannella Publisher Schloss Dagstuhl- Leibniz-Zentrum fur Informatik GmbH, Dagstuhl Publishing ISBN (Electronic) 9783959770767 State Published - 1 Jul 2018 Event 45th International Colloquium on Automata, Languages, and Programming, ICALP 2018 - Prague, Czech Republic Duration: 9 Jul 2018 → 13 Jul 2018 Publication series Name Leibniz International Proceedings in Informatics, LIPIcs Volume 107 ISSN (Print) 1868-8969 Conference 45th International Colloquium on Automata, Languages, and Programming, ICALP 2018 Country/Territory Czech Republic City Prague Period 9/07/18 → 13/07/18 • Analysis and algorithmics • Data reduction rules • Kernelizations • NP-hard problems • Parameterized complexity ASJC Scopus subject areas Dive into the research topics of 'How hard is it to satisfy (Almost) all roommates?'. Together they form a unique fingerprint.
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Curriculum for a Teach-The-Teacher Course? As a starting point, I'd first of all look at teh exisiting curriculum and see if there is something currently covered that scratch might do better. For example, in the UK curiculum in year 3 (7-8 yrs) there is an introduction to control ssytems that, currently, almost everyone currently teaches using "logo" or a derivative therof. That's an obvious candidate for being superceded by Scratch. Usign scratch to replace the exisitng curriculum products can go outside of the ICT class, as well. Your need to have introduced scratch first, but you can show the teachers some of the sample book reports that chn have produced and submitted. I've also seen teaching units where teh chn have to create their own, short, "choose your own adventure" book, somethign that could be translated into scratch very easily. I'm keen - should I start teaching slightly older classes in teh future - to employ a cross curricular approach where (for example) we write a script in English class, draw the characters in art class and then animate the story using scratch in ICT. Finally - let the teachers know that if they themselves start getting interested in scratch, they can use it to produce simple interactive resources to use in teaching, customised to their own I've used it to teach directions (directions game 1 and 2), past tense endings (Ed the cat), RE (the cleansing of the temple) and counting in 10ps (coin counting game), and may also use it in the future in the ICT "introduction to modelling" unit. As soon as I get around to writing a scratch simulation of growing a plant, that is
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Reconstruction of piano hammer force from string velocity A method is presented for reconstructing piano hammer forces through appropriate filtering of the measured string velocity. The filter design is based on the analysis of the pulses generated by the hammer blow and propagating along the string. In the five lowest octaves, the hammer force is reconstructed by considering two waves only: the incoming wave from the hammer and its first reflection at the front end. For the higher notes, four- or eight-wave schemes must be considered. The theory is validated on simulated string velocities by comparing imposed and reconstructed forces. The simulations are based on a nonlinear damped stiff string model previously developed by Chabassier, Chaigne, and Joly [J. Acoust. Soc. Am. 134(1), 648–665 (2013)]. The influence of absorption, dispersion, and amplitude of the string waves on the quality of the reconstruction is discussed. Finally, the method is applied to real piano strings. The measured string velocity is compared to the simulated velocity excited by the reconstructed force, showing a high degree of accuracy. A number of simulations are compared to simulated strings excited by a force derived from measurements of mass and acceleration of the hammer head. One application to an historic piano is also presented. In pianos, precise knowledge of the hammer force in terms of amplitude, duration, and shape is essential since it fully determines the resulting free vibrations of the strings. Therefore, it is not surprising that a vast literature exists on this matter in piano acoustics. For years, one commonly used method for deriving the force consisted in multiplying together the measured mass and acceleration of the hammer head.^1,2 This method yields an acceptable rough estimate of the hammer pulse, though it suffers from several limitations.^3 Among these limitations, the oscillations of the hammer shank were highlighted experimentally by several authors in the past,^4,5 and confirmed by recent simulations.^6 With the emergence of new optical techniques and, in particular, of those using high-speed cameras, new methods are now available for measuring the hammer force with great accuracy.^3 In this paper, an alternative method is proposed for determining the hammer force. This method is based on non-contact measurements of the string velocity from which the hammer force is derived through appropriate filtering based on the properties of the wave propagation along the string. The underlying motivation of this method follows from the observation that the hammer is not always easily accessible on all pianos and, in addition, that sophisticated optical methods cannot always be simply implemented in all situations. This fact has been experienced by the author during experiments on historic pianofortes, for which conservation requirements impose to find fast, portable, and non-intrusive measurements. A similar approach was used in the past for the bowed string.^7 The basic principles of the reconstruction method are presented in Sec. II for a linear undamped nondispersive string (a so-called “ideal” string) excited by an imposed hammer force pulse $fH(t)$ of finite duration τ[H] at a given striking position x[H]. The transverse string velocity $vs(t)$ is measured at a given position x[s] along the string. The initial transient part of $vs(t)$ is represented by the finite sum of elementary waves generated by the hammer, and traveling back and forth along the string. For piano strings, it is shown that 2 to 8 waves are sufficient for representing this initial transient adequately for the complete piano compass, the number of waves to be considered depending on both the fundamental frequency f[1] of the note and force width τ[H]. A classical comb filter is applied to $vs(t)$ in order to reconstruct the hammer force. The appropriate choice of x[s] is imposed by the properties of the filter. In Sec. III, the method is applied to simulated piano strings. The simulations are based on a model of a damped nonlinear stiff string recently developed by Chabassier et al.^8 The physical parameters of the simulated strings are measured on real pianos. The string is excited by an imposed hammer force pulse. This input force has realistic amplitude, time width, and shape for each simulated note. The pertinence of the wave filtering method is validated by comparisons between the input force and the reconstructed force. Three different types of wave filters are used in order to cover the complete range of notes of a piano keyboard. The influence of damping, stiffness, and amplitude of the string motion on the estimation of the hammer force are examined. In Sec. IV, the method is applied to measured piano strings. The string velocity is measured with a calibrated electromagnetic sensor. The reconstructed hammer force serves as the input force for the corresponding simulated string, and the resulting simulated velocity is compared to the measured velocity. In some particular cases, the reconstructed hammer force is compared to the force derived from measurements of hammer head acceleration and mass, and the influence of hammer shank motion is discussed. Finally, an example of application of the method to an historic piano with a Viennese action is presented. A. Ideal piano string model In this section, it is assumed that the motion of the piano string is described by a linear wave equation without damping or stiffness terms. The string is characterized by its tension at rest T[0], length L, and linear mass density μ. The string is rigidly fixed at both ends, agraffe (A) and bridge (B), which are considered as perfectly reflecting boundaries. Only one transverse component of the string motion y(x,t) is considered, in the direction of the initial hammer velocity, which is here assumed to be vertical. In this ideal case, the equations of the struck string are written^9,10 $μ∂2y∂t2−T0∂2y∂x2=fH(t)δ(x−xH), ∀ x∈[0,L],y(x=0,t)=0, y(x=L,t)=0,$ where $fH(t)$ is a hammer force of finite duration τ[H] applied at the striking position x[H], and δ is the Dirac delta function. The transverse speed is $c=T0/μ$, the fundamental frequency is $f1=1 /T1=c/2L$, and $Zc=μT0$ is the characteristic impedance.^9 In pianos, the striking point is situated at a distance roughly equal to $L/10$ from the agraffe.^11 The string velocity $vs(t)$ is measured at position x[s]. Figure 1 shows the case where x[s] is situated between agraffe and hammer. Both situations x[s]<x[H] and x[s]>x[H] will be examined in this section. At position x[H], the hammer blow generates two pulses $vsh(t)=fH(t)/2Zc$ traveling in opposite directions: the left wave toward the agraffe (A), and the right wave toward the bridge (B). These elementary pulses reach the sensor position at successive instants of time t[i] whose expressions are given in the Appendix. In Sec. IIB, an intuitive analysis of the wave propagation on the string is presented, showing that the number of elementary pulses (or waves) to consider in the measured string velocity to reconstruct the hammer force depend on the ratio between the force width τ[H] and the string period T[1]. For pianos, this leads to the elaboration of two-, four- and eight-wave schemes, depending on the played note. In Sec. IIC, it is shown that these three schemes are particular cases of a general formulation of the transfer function between string velocity and hammer force that needs to be inverted in order to recover the force. This inversion imposes precautions in the measurements and, in particular, in the appropriate selection of the measuring point x[s]. B. Principles of the force reconstruction method. Wave analysis 1. Bass and medium range of the piano. Two-wave scheme Reconstructing the hammer force from the string velocity is an inverse problem. Knowing $vs(t)$, the goal is to recover the force pulse $fH(t)$ through the analysis of the wave propagation on the string and the reflection of the pulses at the boundaries (agraffe and bridge). Figure 2 illustrates this point in the typical case of a piano note in the bass (or medium) range, where the string velocity waveform recorded near the agraffe is clearly separated into two distinct parts. In the presented case where x[s]<x[H], the first part of $vs(t)$ is proportional to the sum of the incoming hammer force pulse (1) and its reflection at the agraffe (2), whereas its second part is proportional to the sum of the force pulses reaching the sensor after reflection at the bridge (3) and after further reflection at the agraffe (4). Each reflection at either end is characterized by a change of sign. Using the notations defined in Sec. IIA, we can write for the first part of $vs(t)$, $vs(t)=12Zc[ fH(t−xH−xsc)−fH(t−xH+xsc) ].$ The calculation in the Appendix shows that Eq. (2) can be used to reconstruct the hammer force pulse as long as its width τ[H] is such that This condition is obtained by considering that the force pulse (1) goes to zero before the third pulse reaches the sensor. In other words, Eq. (3) means that the number of elementary waves to consider for the reconstruction is equal to the number of waves reaching the sensor during the duration τ[H] of the force pulse width. Since x[H] usually is on the order of $L/10$, the condition in Eq. (3) corresponds to a force pulse duration slightly smaller than the period T[1] of the string's oscillation. In practice, Fig. 3 shows that this condition is fulfilled for most pianos below the note C5. In this range, the force reconstruction method then consists in inverting the two-wave scheme expressed in Eq. (2). This inversion is achieved by applying a dedicated so-called reconstruction filter to the string velocity, as explained in Sec. IIC. Analysis of the filter properties will show why the velocity sensor must be placed near the agraffe for this two-wave scheme. 2. Upper range of the piano. Four- and eight-wave schemes In pianos, the width of the force pulse decreases less rapidly than the period of the string's oscillation when moving toward the treble range (see Fig. 3). As a consequence, the propagating pulses get closer to one another and several reflections can occur during the pulse width τ[H]. It is shown in the Appendix that, under the condition the initial force pulse goes to zero before the fifth reflected pulse reaches the sensor. Therefore, the string velocity can be represented by a sum of four waves during the initial time interval τ [H]. An example is shown in Fig. 4. The frequency analysis of the corresponding reconstruction filter then shows that the measuring point of the string velocity x[s] now must be close to the bridge in order to recover the force properly (see Sec. IIC). In practice, the condition in Eq. (4) required for this four-wave scheme is rather restrictive. It is fulfilled in a small range of notes for most pianos, usually between the notes C5 and C6 (see Fig. 3). For the remaining two upper octaves of the piano (notes C6–C8), a typical situation is shown in Fig. 5. In the example presented here, it is seen that the first six pulses are involved in the expression of $vs(t)$ over the duration τ[H] of the initial force pulse. Following the same reasoning as previously, we could think of considering six waves over this time interval to reconstruct the hammer force successfully, under the condition that the seventh pulse reaches the sensor after extinction of the first pulse. However, it turns out that such a six-wave scheme is unstable. The idea used here is to consider eight waves (instead of six), because it yields again a stable filter. This procedure does not affect the result of the reconstruction: the output signal obtained is then made of the reconstructed pulse of width τ[H] followed by a varying number of zeroes, depending on the note. Another argument for this eight-wave scheme is that it remains valid under the condition which, in practice, is fulfilled for the remaining upper part of the piano (notes C6 to C8). As for the four-wave scheme, the properties of the corresponding filter imposes to select x[s] close to the bridge (see Sec.IIC). In summary, the force reconstruction is governed by two important criteria. In the time-domain, the width τ[H] of the hammer force pulse imposes the number of waves to consider in the inverse filtering. In the frequency-domain, the properties of the reconstruction filter impose the location of the measuring point x[s] along the string, as shown below in Sec. IIC. C. Reconstruction filters 1. General formulation The string wave propagation qualitatively presented in Sec. IIB can be formally described by a filter with the hammer force as input signal, and the string velocity at the measuring point as output signal. For discrete signals sampled at frequency $fe=1/Te$, the appropriate tool is the z-transform where a delay of T[e] corresponds to a multiplication by $z−1$. Therefore $z−m$ corresponds to a delay of m samples, that is, mT[e] seconds. For details on the z-transform, see, for example, Ref. 12. In what follows, $Vs(z)$ and $FH(z)$ denote the z-transforms of $vs(t)$ and $fH(t)$, respectively. In the Appendix, the example of a plucked string with a pickup output developed by Karjalainen et al. is adapted to the present piano case.^13 For a lossless nondispersive string, it is shown that the transfer function between string velocity and hammer force is given by where n[1] accounts for the propagation delay between hammer and sensor, m[1] and m[2] account for the delays due to hammer and sensor positions, and m[3] is the delay for one complete loop. Conversely, if $Vs(z)$ is known, the hammer force $FH(z)$ can be retrieved through the inverse transfer function If the sensor is situated on the left-hand side of the hammer (on the agraffe side: x[s]<x[H]), it is shown in the Appendix that $n1=(xH−xs)fec; m1=2xsfec=xsfeLf1; m2=(L−xH)feLf1; m3=fef1.$ Conversely, for x[s]>x[H] (sensor on the bridge side), the delay constants are obtained through permutation of x[s] and x[H] $n1=(xs−xH)fec; m1=2xHfec=xHfeLf1; m2=(L−xs)feLf1; m3=fef1.$ 2. Approximate two-, four-, and eight-wave filters The transfer function expressed in Eq. (7) is valid for an infinite number of waves reaching the sensor. The filters corresponding to the wave analysis presented in Sec. IIB are obtained by truncating this expression to the initial portion of the signal. Considering the first two incoming waves only, then H[FV] reduces to ${ HFV2w(z)=2Zczn11−z−m1 for xs<xH,or2Zczn11−z−m2 for xs>xH.$ Notice that the first expression in Eq. (10) is the transfer function directly obtained by taking the z-transform of Eq. (2). Similarly, a four-wave approximation is given by Finally, for the eight-wave scheme, the corresponding filter is obtained by approximating the factor $(1−z−m3)$ by the two first terms of its expansion, which yields Notice, in this case, that Eq. (7) can be used alternatively instead of Eq. (12). The advantage is then to deal with the feed-forward filter $1−z−m3$ rather than with the feedback version $1/(1+z−m3) $, which adds high peaks in the spectrum. After filtering of the velocity, the hammer force pulses are obtained by truncating the signal in time below the upper limit of validity for each of the three schemes given in Eqs. (3)–(5). Using the general form of the inverse filter Eq. (7) theoretically yields an infinite domain of validity in time. However, in practice, due to the presence of noise and progressive departure from the underlying model when dealing with real signals, the filtered velocity again has to be truncated to reasonable values. An upper bound of 10ms is, in general, sufficient for recovering the force. 3. Selection of the measurement point x[s] Some precautions should be taken in using the previous filters to reconstruct the hammer force. This is due to the fact that both delays m[1] and m[2] now introduce comb filtering effects with poles at frequencies $fkm1=kfe/m1$ and $fkm2=kfe/m2$, respectively. Because of their high gain, these pole frequencies can perturb the reconstruction of the force if they are smaller than the upper limit f[VM] of the velocity spectrum. These perturbations are due to amplification of noise and deviations from the ideal case in real signals at these frequencies. Therefore, given the expressions of m[1] and m[2] in Eqs. (8) and (9), the appropriate strategy for limiting these perturbations is to select x[s] so that f[VM] is smaller than the lowest poles $fe/m1$ and $fe/m2$, which yields the following conditions: ${ for xs<xH: xsL<f1fVM and L−xHL<f1fVM,for xs>xH: xHL<f1fVM and L−xsL<f1fVM.$ For x[s]<x[H], the second condition, relative to m[2] in Eq. (13) is unacceptable since it implies that the spectrum of the string velocity should be limited to the fundamental frequency. Therefore, only the condition on m[1] remains, which means that this configuration is valid for the two-wave scheme only. According to the wave analysis presented in Sec. IIB, this scheme is applicable to the bass and medium range of the piano. In this range, x[s] must be small compared to L (sensor close to the agraffe) so that the poles of the filter are rejected out of the velocity spectrum. With $xs/L=100$, for example, the lowest pole of the filter is rejected beyond the 100th harmonic of the string. The configuration x[s]>x[H] is applicable to four- and eight-wave schemes if the two corresponding conditions in Eq. (13) are fulfilled. From the wave analysis in Sec. IIB, we know that such schemes are applicable in the treble range between C5 and C8. The second condition on m[2] implies that x[s] is close to L, which is obtained now by placing the sensor close to the bridge. In practice, good results are obtained for $0.05<(L−xs)/L<0.1$. The first condition related to $xH/L$ is imposed by the construction of the piano. On average $xH/L=0.1$, which means that the lowest pole frequency of the reconstruction filter is rejected beyond the tenth harmonic. Usual piano tones show relatively little energy beyond this harmonic, between C5 and C8 (see, for example, Fig. 11), and thus the configuration x[s]>x[H] is valid in this range. Notice also that the regularization presented in Sec. IIC4 allows to reduce the perturbating effects of the poles. 4. Further considerations: Dispersion, fractional delays and regularization In the bass range of the piano, it becomes necessary to account for the dispersion of the wave due to string stiffness. One widely used method for this is to divide the transfer function $HVF(z)$ in Eq. (6) by an appropriate allpass filter D(z) whose phase and delay properties are based on the physical parameters of the string. For a stiff string, the phase shift $φ(f)$ for a roundtrip is given where $B=π2EI/T0L2$ is the stiffness coefficient. The associated group delay $γ(f)$ is given by $γ(f)=1f11[ 1+Bf2f12 ]3/2.$ The pertinent dimensionless coefficient here is $εs=Bf2/f12$ which, for a given frequency f, continuously decreases from bass to treble. For reproducing this delay, a set of eight biquad filters in cascade was applied, following the method proposed by Abel et al.^14,15 The delays defined in Eqs. (8) and (9) are usually not integers. Therefore one useful strategy is to use fractional delay filters to represent them. Here, Thiran allpass filters were tested successfully.^16 Oversampling is another strategy, so that approximating n[1], m[1], m[2], and m[3] by the nearest integers yields a small error. Both methods gave comparable results. Finally, a regularization can be made in replacing $(1−z−m1)$ and $(1−z−m2)$ by $(1−g1z−m1)$ and $(1−g2z−m2)$ in Eq. (7), with $0.99<g1,g2<0.999$, so that the gain of $HFV(z)$ does not tend to infinity at those frequencies. The consequences on the pole frequencies are negligible. Physically, such a regularization accounts for a slight dissipation at both ends of the string. The interval of values of g[1] and g[2] are selected somewhat arbitrarily, but the examples shown in Secs. III and IV show that this has no appreciable effect on the amplitude of the reconstructed force. However, selecting g[1] and g[2] smaller than 0.99 (which would correspond to losses at the ends higher than 1%) generally leads to unwanted drift in the reconstructed force waveform, showing that the corresponding damping at the string ends is then probably overestimated. Notice that, for coherence, the factor $(1−z−m3)$ must be then replaced by $(1−g1g2z−m3)$. In conclusion, after inversion, the most general formulation of the transfer function used for reconstructing the hammer force is written where the delays are implemented with fractional delay filters. In this section, the method of hammer force reconstruction presented in Sec. II is tested on simulated piano strings. Using string simulations allows a perfect control on the input hammer force and careful comparison with the reconstructed force. In addition, the influence of the main causes of departure from the ideal string case on the quality of the reconstruction can be tested and quantified separately, which is not possible on real strings. Specifically, we take benefit of the simulation model to test the influence of absorption, dispersion, and non-linearity of the string wave propagation on the resulting hammer force reconstruction. The waves schemes presented in Sec. II (two-, four- and eight-waves, respectively) are tested in the different note ranges of the keyboard of various pianos: bass, medium, and treble, and for different striking forces (p,mf,f). The reconstructed hammer force is obtained by filtering the simulated string velocity at a given sensor position x[s]. A. Model and method The string model used for the simulations is a part of a complete numerical piano model developed by Chabassier et al.^8,17 This model will not be detailed further here, since it has been presented in these previous papers extensively. For the string part, this model includes damping, stiffness, and nonlinear terms. The frequency-dependent damping is classically modeled by two-terms: a so-called “fluid” term and a “viscoelastic” damping term for each component of the string motion. The stiffness term is obtained by modeling the string as a prestressed beam. The additional nonlinear terms account for the influence of the amplitude on the string motion.^18 The transverse motion of the string is assumed here to be vertical (in the direction of the hammer force) and coupled to its longitudinal motion. The motion of both ends (agraffe and bridge) is ignored in the description used here. The values of the parameters are derived from measurements on real strings (see Table I). TABLE I. Note . D♯1 . D♯1 . F♯2 . C4 . C♯5 . C♯5 . D♯6 . E6 . Piano JBS73 JBS50 JBS73 JBS73 JBS36 JBS73 BSD BSD L (m) 1.717 1.710 1.57 0.65 0.261 0.335 0.165 0.137 $Φ$ (mm) 1.25 1.1 1.14 0.94 0.73 0.885 0.80 0.80 T[0] (N) 781.6 520 598 552.6 267.9 593 680.5 522.8 f[1] (Hz) 36 37 87 245 547 523 1258 1328 m[H] (g) 5.9 NM 5.4 4.8 NM 4.2 7.02 6.96 x[H] (mm) 190 180 165 75 28 42 20 17 $xs$ (mm) 10 9 27 15 14 325 156 127 $μ=ρA$ (gm^−1) 51.1 32.5 8.01 5.45 3.22 4.83 3.95 3.95 $A=πΦ24, I=πΦ464, E=2.0×1011 Pa$ Note . D♯1 . D♯1 . F♯2 . C4 . C♯5 . C♯5 . D♯6 . E6 . Piano JBS73 JBS50 JBS73 JBS73 JBS36 JBS73 BSD BSD L (m) 1.717 1.710 1.57 0.65 0.261 0.335 0.165 0.137 $Φ$ (mm) 1.25 1.1 1.14 0.94 0.73 0.885 0.80 0.80 T[0] (N) 781.6 520 598 552.6 267.9 593 680.5 522.8 f[1] (Hz) 36 37 87 245 547 523 1258 1328 m[H] (g) 5.9 NM 5.4 4.8 NM 4.2 7.02 6.96 x[H] (mm) 190 180 165 75 28 42 20 17 $xs$ (mm) 10 9 27 15 14 325 156 127 $μ=ρA$ (gm^−1) 51.1 32.5 8.01 5.45 3.22 4.83 3.95 3.95 $A=πΦ24, I=πΦ464, E=2.0×1011 Pa$ In the simulations presented in this section, the hammer forces are given as input data. The amplitude, duration, and shape of the imposed force pulses are selected as close as possible to the reality. In some cases, the hammer mass and acceleration were measured on real strings beforehand and, thus, the input force is obtained by multiplying together these two quantities. In the other situations, the string-hammer coupling is modeled by a standard power law in which the coefficients are adjusted in order to obtain realistic hammer force amplitude and duration.^19,20 As it has been already pointed out by several authors,^3,6 and as discussed further in Sec. IV of this paper, we are aware of the fact that these input hammer pulses might differ from the real forces, either because of several causes of experimental errors (hammer misalignment, vibrations of the hammer shank,…), or because of the approximate modeling by a power law.^5 However, these differences do not invalidate the testing of the reconstruction procedure presented here. In fact, this reconstruction is supposed to account for any hammer pulse, provided that the linear theory presented in Sec. II is applicable, and for the appropriate wave filter (two-waves, four-waves, or eight-waves) compatible with the pulse duration τ[H]. In Sec. IIIB, it is analyzed to what extent realistic departures from the “ideal linear string,” in terms of absorption, dispersion, and nonlinear propagation, can affect the reconstruction of the hammer force. A systematic study of the influence of each of these perturbation terms is presented. In Sec. IIIC, representative examples of hammer force reconstruction are presented in the three main ranges of the piano: bass, medium, and treble, using the appropriate waves schemes. B. Effects of some causes of departure from the ideal string on the reconstructed force 1. Absorption and damping Consider a string with fundamental frequency $f1=1/T1$ and frequency-dependent time-constant $τ(f)$ due to damping. For each partial n, the attenuation due to damping during a roundtrip is proportional to $1−exp[−T1/τ(fn)]$, which, for $T1≪τ$, reduces to $T1/τ$. An order of magnitude can be obtained by calculating this attenuation around the mid-frequency $fm$=1kHz. As an example, we measured $τ(fm)$ =1.2, 0.7, and 0.2s, for the three strings D♯1, C4, and E6 examined here. This yields an attenuation of 2.3%, 0.6%, and 0.4%, respectively, for one roundtrip on these three strings. In practice, the observed attenuation is smaller than these limits (<1%) since the distance traveled by the wave is less than twice the string length (see Fig. 6). As a consequence, no correction was applied here since the attenuation error is weak compared to the other sources of error. Notice however that, if necessary, an additional filter can easily be designed for compensating the frequency-dependent damping along the string.^21 2. Dispersion due to stiffness At f=1kHz, the stiffness coefficient $εs$ defined in Sec. IIC takes the values 7.9×10^−2 for the string D♯1, 1.28×10^−2 for the string C4, and 2.3×10^−3 for the string E6. This results in group delays (normalized to the period) equal to 0.12, 1.9×10^–2, and 3.4×10^–3, respectively, for strings D♯1, C4, and E6. As a consequence, dispersive effects are clearly seen on velocity waveforms for the strings in the bass range, which results in a distortion of the wave, even for reduced propagation distances on the order of x[H] (see Fig. 6). For the two lowest octaves of most pianos, it was thus necessary to design allpass filters in order to compensate for this dispersion (see Fig. 12). The dispersive effects decrease rapidly with the key number, and do not alter the reconstruction of the force for the notes above A2 in all observed pianos. 3. Amplitude non-linearity An appropriate strategy for evaluating the influence the non-linearity due to the amplitude of the string motion is to examine the nonlinear vibrating equation, which, to a first-order expansion, can be written as^22 $∂2y∂t2≈c2∂2y∂x2[ 1+32cL2c2(∂y∂x)2 ],$ where $cL2=EA/μ$ is the longitudinal wave speed in the string of cross-sectional area A and linear density μ. An order of magnitude for the derivative $∂y/∂x$ is given by the ratio V/c between the amplitude of the string velocity V and the transverse speed c. In total, the relevant dimensionless non-linearity coefficient can be written This coefficient quantifies the change in transverse speed due to the amplitude of the wave, which might induce some distortion. Here, again, $εNL$ regularly decreases as the pitch of the note (or, equivalently, the key number) increases. For V=1m/s (which, in most cases, corresponds to a mezzoforte level), $εNL$ takes the values 2.39×10^–2, 3.7×10^–3, and 2.83×10^–3 for the three strings D♯1, C4, and E6, respectively. In practice, the effects of non-linearity on the force reconstruction were observed in the low bass range only, for V higher than 1m/s (see Fig. 6). Notice that Eq. (17) accounts for the transverse effects of the amplitude non-linearity only, which is equivalent to a modulation of tension. The general nonlinear model used here for the simulations also accounts for the longitudinal component of the string, which is not measured by the sensor. In summary, the above dimensionless analysis of the three “perturbation” terms (damping, stiffness, amplitude non-linearity) in the ideal string wave equation shows that stiffness is likely to be the main cause of errors in the force reconstruction. Nonlinear and damping effects are much weaker, although the errors due to non-linearity might rapidly increase with amplitude, since they are proportional to the square of the string velocity. Another important result of this analysis is that the effects of these perturbation terms regularly decrease with the fundamental frequency of the note. It is therefore expected that they are dominant in the bass range. The consequences on the reconstruction of the hammer force are shown in Sec. IIIC for some simulated notes in the bass, medium, and treble ranges of the piano. C. Examples of reconstructed hammer forces 1. Bass range. Two-wave scheme Figure 6 shows an example of reconstructed hammer force for the string D♯1 of a pianoforte made by J. B. Streicher in 1873 (JBS73). The excitation corresponds to a level between forte and fortissimo, with a maximum force of 32N, and a maximum string velocity of V=2.5m/s. A two-wave scheme is applied. Figure 6(a) shows that the reconstructed force perfectly coincides with the input force for a simulated ideal string. The errors (in N) due to damping and nonlinear terms are shown in Fig. 6(b). The maximum discrepancy due to the damping terms is 0.02N, with corresponds to a relative error of 6.0×10^–4, and can be ignored. The maximum error due to the presence of nonlinear terms in the string wave equation is equal to 0.36 N, or 1.13% in relative value. This error is comparable to other potential causes of errors encountered in the application of the method to real strings (see Sec. IV). Finally, Fig. 6(c) shows the effects of stiffness on the reconstructed force. In this case, the consequences of the dispersion are clearly seen, resulting in an underestimation of the maximum force (around 20% here), a slight broadening of the pulse, and the presence of characteristic oscillations in the onset. The results presented here are representative of those observed in the two lowest octaves of the investigated pianos. The discrepancies due to stiffness and non-linearity increase when approaching the lowest notes of the keyboard. They decrease when moving to the higher keys and, in general, as the amplitude of the hammer force decreases. 2. Medium range. Two-wave scheme In this paragraph, representative examples of hammer force reconstruction in the medium range of the piano are presented. This range covers nearly three octaves (from note E2 to C5). In this interval of notes, a two-wave reconstruction scheme is applied. The string parameters and the input force pulses are derived from measurements performed on a JBS73 piano (see Table I). In Fig. 7, the input forces (solid lines) and the reconstructed forces (dashed lines) are shown for the notes F♯2 (fundamental 87Hz) and C4 (fundamental 245Hz). The reconstruction is made with the complete string model, including damping, stiffness, and nonlinear terms. It can be seen that the hammer forces are very well reproduced in each case. The amplitudes, time duration, and shapes of the reconstructed forces coincide with those of the input forces almost exactly. This means that the different causes of departure from the “ideal” string integrated in the simulation model do not affect the wave propagation substantially. As a consequence, the effects of the additional terms (damping, stiffness, non-linearity) remain small. However, one exception should be underlined: as seen in Fig. 7, the first pulse of the reconstructed hammer force for the note F♯2 shows a slightly more rounded profile than the input force, which results in a smaller amplitude by a few percent. This discrepancy is due to the dispersive effects of stiffness. 3. Treble range. Four- and eight-wave schemes Finally, reconstructed forces are presented here for the notes in the upper part of the keyboard. This range corresponds to the domain of notes where the period becomes smaller than the width of the hammer pulse, as shown in Fig. 3. In this domain (with pitch higher than C5 for most pianos), four- and eight-wave schemes must be applied for reconstructing the hammer force. Figure 8 shows two examples of reconstruction with the complete string model. The hammer force of the C♯5 string is reconstructed with a four-wave scheme, while the E6 force is reconstructed with an eight-wave scheme. In both cases, no differences can be seen between input and reconstructed forces, which means that the perturbation terms have no visible effects. In conclusion of this section, it is confirmed that the hammer forces can be reconstructed from simulated string velocities generated by a model that includes damping, stiffness, and nonlinear terms in addition to the main inertial and tension terms. Depending on the notes, two-, four-, or eight-wave filters must be applied to the string velocity to recover the hammer force. As shown in Sec. II, the choice of the filter depends on the ratio between the period of the string oscillation and the duration of the hammer pulse. In the low bass range, corresponding to the two lowest octaves for most pianos, dispersion effects due to stiffness are more pronounced, which alters the quality of the reconstruction. In this case, the dispersion needs to be corrected by the use of an allpass filter whose parameters are derived from the physical parameters of the string, as seen in Sec. IIC. Due to the relatively small distance traveled by the waves along the string during its transient oscillation, the other causes of departure from the “ideal” “string” (damping, non-linearity) induce only small effects in the reconstruction of the hammer force: the duration and shape of the pulse are not affected, and its magnitude is estimated within a few percent, or less, compared to the input After validation on simulated string waveforms, this section shows some results obtained with the hammer force reconstruction method applied to real piano strings. In addition to damping, stiffness and non-linearities, other physical phenomena, might alter the quality of the force reconstruction. The coupling between the strings of a given note, the whirling motion of the string, and lack of accuracy in the velocity measurements, for example, might lead to wrong estimations. Therefore, it is challenging to test whether the proposed method is robust enough against these potential causes of errors. The general procedure is shown in Fig. 9. In the most general case (case 1), the hammer force is not known. The hammer force $fHR(t)$ is then reconstructed through filtering of the measured string velocity $vsm(t)$, following the method presented in Sec. II. This force is used as input data in the string simulation model presented in Sec. III, which yields the simulated string velocity $vsi (t)$, for further comparisons with $vsm(t)$, as shown in Sec. IVB. In the present study, an estimation of the measured hammer force $fHM(t)$ was also derived in some particular cases from measurements of mass and acceleration of the hammer head. In these cases (case 2), this measured hammer force is compared to the reconstructed force, and also serves as input data in the string simulation model, in order to compare the resulting string velocity $vSM(t)$ to both $vsm(t)$ and $vsi(t)$. This second procedure is detailed in Sec. IVC. The section ends with an example of application of the method to the case of an historic piano with a Viennese action. A. Measurements Figure 10 shows a picture of the setup used for measuring the string velocity. The sensor is a standard electromagnetic transducer made of a coil wound onto a small cylindrical permanent magnet, similar to an electric guitar pickup, and placed in the vicinity of the string. This transducer delivers a voltage proportional to the vertical component of the string velocity. The distance between the magnet and the string is adjusted with a micrometer screw. The sensor and the screw can slide from string to string along a horizontal metallic bar placed over the instrument. A calibration procedure is made prior to the measurements, where the signal delivered by the transducer is compared to the output of a laser vibrometer Polytec IVS-400 (Hörsching, Austria). The efficiency of the transducer (in mV/m/s) is tested against string amplitude, frequency, and sensor-string distance. The efficiency is constant for a string amplitude between 20 and 300μm, and shows a slight increase of 5% between 0.3 and 1.0mm. In our experiments, the sensor is placed either close to the bridge or to the agraffe, so that the string amplitude is <200μm. B. Comparison between measured and simulated string velocity We show here some representative results of the force reconstruction method applied to measured string velocities. This is supposed to correspond to the standard, and most useful, application of the method. In the absence of measured hammer force, the validity of the reconstructed force is tested by comparing the measured string velocity with simulations where the input signal is the reconstructed force (see Fig. 9). Figure 11 shows the results of the hammer force reconstruction method applied to the string C♯5 of a J. B. Streicher piano built in 1836 (JBS36). It can be seen that the string velocity simulated with the reconstructed force reproduces the measured waveform almost perfectly, even on a large time scale. This is confirmed by the spectral analysis performed on the first 50ms of the signal, which shows a high degree of coincidence between 0 and 5kHz. A DC-component at −30dB below the maximum can be seen on the measured spectrum. In the bass range, an allpass filter is applied to the measured string velocity in order to compensate for the dispersion due to stiffness. Figure 12(b) shows the effect of this filter on the hammer pulse when the dispersion is reduced, which is comparable to the results obtained on the simulations in Sec. IIIC when the stiffness term is removed from the string equation. The string velocity simulated with the filtered reconstructed force shows a good agreement with the measured velocity over a large time scale (35ms) despite some discrepancies in the amplitudes of some sub-oscillations. These discrepancies could be due to some inaccuracy in the modeling of the frequency-dependent damping, always a delicate task, especially in the bass range. However, the comparison between the two corresponding spectra shows an excellent agreement. The major differences are seen at the peaks of small amplitudes, around -20dB below the maximum (0dB), and can be attributed either to some small errors in the estimation of the string parameters (damping) or in the measurements (precise location of the string sensor). C. Comparison between reconstructed and measured force Figure 13 shows an example of hammer force reconstructed from measurements of the string velocity for the note F♯2 of a J. B. Streicher piano (1873). The reconstructed force is compared to the measured force derived from measurements of mass and acceleration of the hammer head. One can see that measured and reconstructed forces coincide during the first pulse of the waveform. However, some differences exist in the magnitude of the second pulse. In addition, the measured force tends faster to zero and the remaining oscillations of the hammer head are clearly visible. Both forces are used as input forces for simulating the string velocities, using the measured parameters shown in Table I. The simulated string velocities are then compared to the measured velocity. Here, again, one can see an excellent agreement between the measured velocity and the velocity simulated with the reconstructed force, even on a long period of time. The simulations performed with the reconstructed force as input quantity yields waveforms and spectra closer to the measurements than the simulations using measured hammer head acceleration and mass at the input. This is not surprising since the reconstruction filter does not take the vibrations of the hammer shank into account. In the treble range, Fig. 14 shows the results of the hammer force reconstruction method applied to the string D♯6 of a Bösendorfer piano with an eight-wave filtering scheme. The comparison between the reconstructed force and the measured force obtained through multiplication of mass and acceleration of the hammer head again shows that this latter is affected by the oscillations of the head + shank system, resulting in differences in duration and shape of the main pulse. Again, and for the same reasons as previously, the string velocity simulated with the reconstructed force shows a better agreement with the measured velocity that the velocity simulated with the measured force truncated to its first positive pulse. This agreement is confirmed by observing the velocities over a larger period of time (8ms) and on the spectra, where only small differences are visible below 5kHz. In summary, comparing the hammer pulses derived from the present method with hammer head mass-acceleration measurements shows that the present method is insensitive to the perturbing oscillations of the hammer shank. Its main advantage also follows from the fact that it does not need to have direct access to the hammer: this feature is of the utmost importance when measuring valuable and fragile instruments such as historic pianofortes. D. Application to historic pianofortes One running application of the present method is to investigate the evolution of hammer-string interaction in the history of piano making. In this respect, Fig. 15–(a) shows an example of a reconstructed hammer force (note D3) for a copy recently made by G. Hecher of a pianoforte built by Nanette Streicher in 1805 (GH05). This instrument uses a Viennese action,^23 and 200-yr-old leather is glued on the hammer heads. From the qualitative comparison with the hammer acceleration measured of a modern piano in the same half-octave (Steinway D, note G3) in Fig. 15(b), one can see that both shapes are globally similar. However, close observation shows some differences in the slopes and rounded profiles of the pulses, which are likely to be attributable to the respective behavior of leather and felt. Also interesting is the comparison between measured and simulated velocity for the GH05-D3 string in Fig. 15(c). The first part of both waveforms are identical, but significant differences now exist in the second part where the measured oscillations are significantly damped compared to the velocity simulated with the reconstructed force at the input. A very plausible explanation for this difference lies in the specificity of the Viennese action, where the hammer slips along the string before leaving it, a feature that is not taken into account in the reconstruction model. In this paper, an original method was presented for reconstructing the piano hammer force from measurements of the string velocity. This method is based on the use of dedicated filters applied to the measured velocity. The design of these filters is derived from the analysis of wave propagation along a lossless linear string. For the five lowest octaves of a piano, a filter based on two-wave propagation (one incident and one reflected wave) is sufficient for reconstructing the hammer force. For the highest notes, a four- or an eight-wave scheme is necessary. The number of waves to consider is imposed by the duration of the hammer pulse. The velocity sensor must be placed either near the agraffe (two-wave scheme) or near the bridge (four- and eight-wave schemes) in order to ensure that the reconstructed force is not perturbed by the poles of the filter. In a first step, the force reconstruction method has been validated on simulated piano strings, where the model includes absorbing, dispersive, and nonlinear terms. The input parameters of these simulations were derived from measurements on a large variety of historic and modern pianos. The results show that both the damping and nonlinear terms only have small effects on the estimation of the hammer force. On the contrary, the dispersion due to stiffness leads to significant discrepancies in the two lowest octaves, which needs to be compensated by means of a suitable allpass filter. In a second step, the force reconstruction method is tested on velocity measurements performed on real piano strings. The measured string velocity is compared to the corresponding simulated velocity excited by the reconstructed force. It is found that the reconstructed hammer force yields a simulated string velocity that reproduces the measured velocity with a high degree of accuracy, both in time and frequency. This proves that the simplified string model used for designing the reconstruction filters accounts well for the observed wave propagation during the first oscillations after the hammer stroke. During this small time interval, known complicating effects, such as the whirling motion of the strings and the coupling between adjacent strings of the same note, are not yet initiated and do not perturb the reconstruction of the force. In this context, it would be of interest to compare the results with methods requiring the use of longer portions of velocity signal as, for example, those based on the use of direct and inverse Fourier transforms in the frequency domain. Comparisons of the reconstructed force with the results obtained by other experimental methods (such as optical methods and high-speed cameras) also remain to be done in order to completely assess the pertinence of the present filtering. In some cases, the reconstructed force was compared to an estimate of the hammer force obtained by multiplying together the measured mass and acceleration of the hammer head. These comparison shows reasonable similarities in terms of amplitude, shape, and duration, but differences are observed, which are mainly attributable to the oscillations of the hammer shank. Simulations of the string velocity performed with this measured force at the input of the present filters show more discrepancies than with the reconstructed force when compared to the measured velocity: such a result is not surprising if we consider that the dynamics of the shank are not included in the underlying model on which the reconstruction is based. In this context, extending the method to a filter with a transfer function between hammer acceleration and string velocity, which would include a model of shank, could be an interesting perspective. Finally, the reconstruction method is applied to a copy of an historic piano equipped with a particular action (the so-called Viennese action). Qualitative comparison between the reconstructed force and the hammer acceleration measured on a neighboring note on a modern piano shows some differences that might be due to distinct compression behavior of leather and felt, respectively. The comparison between measured and simulated string velocity of the D3 note played on the historic piano also shows interesting features that are likely attributable to the specificity of the Viennese action, where the hammer slips along the string before leaving it. This is known to produce a softer sound than modern pianos. This project was supported by a Lise-Meitner-Fellowship of the Austrian Science Fund (FWF; Project number M 1653-N30). The author wishes to thank Alex Mayer (University of Music and Performing Arts Vienna), Caroline Haas and Michael Kirchweger (Technical Museum Vienna), and Gert Hecher (Das Klavier-Atelier, Vienna) for their valuable help in the measurements. The simulations presented in this paper were carried out using the Plateforme Fédérative pour la Recherche en Informatique et Mathématiques (PLAFRIM) experimental platform, being developed under the Inria PlaFRIM development action with support from Institut Polytechnique de Bordeaux, Laboratoire Bordelais de Recherche en Informatique, and Institut de Mathématiques de Bordeaux and other entities: Conseil Régional d'Aquitaine, Université de Bordeaux, and Centre National de la Recherche Scientifique (and Agence Nationale de la Recherche) in accordance with the programme d'investissements d'avenir (see http://www.plafrim.fr/ ). The author is indebted to Juliette Chabassier and Marc Duruflé for their assistance in using the PlaFRIM environment. 1. Wave analysis For x[s]<x[H], the pulses generated by the hammer reach the sensor at the successive instants of time $t1=(xH−xs)/c; t2=(xH+xs)/c; t3=(2L−xH−xs)/c; t4=(2L−xH+xs)/c.$ For a hammer pulse of width τ[H], the condition for non-overlapping with the third pulse within the time interval τ[H] is given by Under this condition, the velocity signal is the sum of two waves between t[1] and $t1+τH$. For x[s]>x[H] and $L−xs<xH$ (sensor close to the bridge), we have $t1=(xs−xH)/c; t2=(2L−xH−xs)/c; t3=(xH+xs)/c; t4=(2L+xH−xs)/c.$ The subsequent arrival times are such that $ti+4=ti+2L/c$. The condition of non-overlapping with the fifth pulse between t[1] and τ[H] then becomes Under this condition, the hammer force can be reconstructed with the velocity represented by four waves. Finally, the condition of non-overlapping with the ninth pulse within the interval τ[H] is For $T1<τH<2T1$, the force can then be reconstructed with the velocity signal represented by eight waves. 2. Reconstruction filter The transfer function between string velocity and hammer force for a lossless nondispersive struck string expressed in Eq. (6) can be conveniently obtained with the help of the dual delay-line model shown in Fig. 16, adapted from the plucked string model with pickup by Karjalainen et al.^13 The upper delay line accounts for the wave propagation in the direction from agraffe (A1) to bridge (B1), through hammer (H1) and sensor (S1) positions. The lower delay line accounts for the propagation from bridge (B2) to agraffe (A2). The hammer force F[H] gives rise to two pulses $FH1$ and $FH2$ at position x[H], and the string velocity at position x[s] is the sum of the contributions $Vs1$ and $Vs2$ from both lines. In the case x[H]<x[s], the delays (in samples) between the various points $nH=xHfe/c; ns−nH=(xs−xH)fe/c; nL−ns=(L−xs)fe/c.$ R[a] and R[b] are the reflection coefficients at agraffe and bridge, respectively. On the right-hand side of the lines, we can write $Vs=Vs1+Vs2=Vs1[ 1+Rb z−2(nL−ns) ].$ Denoting $VH1$ the string velocity at the hammer position traveling to the right, we can write at the sensor position $Vs1=VH1 z−(ns−nH)+RaRb z−2nL Vs1.$ Finally, at the hammer position, we can write $VH1=12FHZc[ 1+Ra z−2nH ].$ From Eqs. (A7)–(A9), we get $HVF=VsFH=12Zcz−n1[ 1+Ra z−m1 ][ 1+Rb z−m2 ]1−RaRb z−m3,$ with $n1=ns−nH, m1=2nH, m2=2(nL−ns)$ and $m3=2nL$. For perfect reflecting ends, we have $Ra=Rb=−1$. For dissipative ends, we can write $Ra=−g1$ and $Rb=−g2$, with $| g1 |,| g2 |<1$. , “ Model for piano hammers: Experimental determination and digital simulation J. Acoust. Soc. Am. D. E. , “ Piano string excitation V: Spectra for real hammers and strings J. Acoust. Soc. Am. , “ Experimental investigation of the piano hammer-string interaction J. Acoust. Soc. Am. E. V. , “ From touch to string vibrations. II: The motion of the key and hammer J. Acoust. Soc. Am. J. P. , “ Piano hammers and their force compression characteristics: Does a power law make sense? J. Acoust. Soc. Am. , “ Energy based simulation of a Timoshenko beam in non-forced rotation. Influence of the piano hammer shank flexibility on the sound J. Sound Vib. R. T. , and , “ Reconstruction of bowing point friction force in a bowed string J. Acoust. Soc. Am. , and , “ Modeling and simulation of a grand piano J. Acoust. Soc. Am. P. M. Vibration and Sound Acoustical Society of America Melville, NY ), Chap. 3, pp. , “ Numerical simulations of piano strings. I. A physical model for a struck string using finite difference methods J. Acoust. Soc. Am. H. A. , “ Design and tone in the mechanoacoustic piano. Part III. Piano strings and scale design J. Acoust. Soc. Am. J. O. Introduction to Digital Filters with Audio Applications W3K Publishing Palo Alto, CA ), Chap. 3, pp. . Available at (Last viewed 3/17/2016). , and , “ Plucked-string models: From the Karplus-Strong algorithm to digital waveguides and beyond Comput. Music J. J. S. , and J. O. , “ Robust, efficient design of allpass filters for dispersive string sound synthesis IEEE Signal Proc. Lett. J. S. J. O. , “ Robust design of very high-order allpass dispersion filters ,” in Proceedings of the 9th Int. Conference on Digital Audio Effects , Montreal ( ), pp. , “ Recursive digital filters with maximally flat group delay IEEE Trans. Circuit Theory , and , “ Time domain simulation of a piano. Part 1: Model description ESAIM: Math. Modell. Numer. Anal. P. M. K. U. Theoretical Acoustics New York ), Chap. , “ Acoustics of pianos Appl. Acoust. , “ Experimental and computational studies of piano hammers Acta Acust. Acust. J. O. , “ Physical modeling using digital waveguides Comput. Music J. G. V. , “ Large amplitude damped free vibration of a stretched string J. Acoust. Soc. Am. , “ The check in some early pianos and the development of piano technique around the turn of the 18th century Early Music
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Automorphisms of Manifolds and Algebraic $K$-Theory: Part IIIsearch Item Successfully Added to Cart An error was encountered while trying to add the item to the cart. Please try again. Please make all selections above before adding to cart Automorphisms of Manifolds and Algebraic $K$-Theory: Part III eBook ISBN: 978-1-4704-1720-8 Product Code: MEMO/231/1084.E List Price: $71.00 MAA Member Price: $63.90 AMS Member Price: $42.60 Click above image for expanded view Automorphisms of Manifolds and Algebraic $K$-Theory: Part III eBook ISBN: 978-1-4704-1720-8 Product Code: MEMO/231/1084.E List Price: $71.00 MAA Member Price: $63.90 AMS Member Price: $42.60 • Memoirs of the American Mathematical Society Volume: 231; 2014; 110 pp MSC: Primary 57; Secondary 19 The structure space \(\mathcal{S}(M)\) of a closed topological \(m\)-manifold \(M\) classifies bundles whose fibers are closed \(m\)-manifolds equipped with a homotopy equivalence to \(M\). The authors construct a highly connected map from \(\mathcal{S}(M)\) to a concoction of algebraic \(L\)-theory and algebraic \(K\)-theory spaces associated with \(M\). The construction refines the well-known surgery theoretic analysis of the block structure space of \(M\) in terms of \(L\)-theory. □ Chapters □ 1. Introduction □ 2. Outline of proof □ 3. Visible $L$-theory revisited □ 4. The hyperquadratic $L$–theory of a point □ 5. Excision and restriction in controlled $L$–theory □ 6. Control and visible $L$-theory □ 7. Control, stabilization and change of decoration □ 8. Spherical fibrations and twisted duality □ 9. Homotopy invariant characteristics and signatures □ 10. Excisive characteristics and signatures □ 11. Algebraic approximations to structure spaces: Set-up □ 12. Algebraic approximations to structure spaces: Constructions □ 13. Algebraic models for structure spaces: Proofs □ A. Homeomorphism groups of some stratified spaces □ B. Controlled homeomorphism groups □ C. $K$-theory of pairs and diagrams □ D. Corrections and Elaborations • Permission – for use of book, eBook, or Journal content • Book Details • Table of Contents • Requests Volume: 231; 2014; 110 pp MSC: Primary 57; Secondary 19 The structure space \(\mathcal{S}(M)\) of a closed topological \(m\)-manifold \(M\) classifies bundles whose fibers are closed \(m\)-manifolds equipped with a homotopy equivalence to \(M\). The authors construct a highly connected map from \(\mathcal{S}(M)\) to a concoction of algebraic \(L\)-theory and algebraic \(K\)-theory spaces associated with \(M\). The construction refines the well-known surgery theoretic analysis of the block structure space of \(M\) in terms of \(L\)-theory. • Chapters • 1. Introduction • 2. Outline of proof • 3. Visible $L$-theory revisited • 4. The hyperquadratic $L$–theory of a point • 5. Excision and restriction in controlled $L$–theory • 6. Control and visible $L$-theory • 7. Control, stabilization and change of decoration • 8. Spherical fibrations and twisted duality • 9. Homotopy invariant characteristics and signatures • 10. Excisive characteristics and signatures • 11. Algebraic approximations to structure spaces: Set-up • 12. Algebraic approximations to structure spaces: Constructions • 13. Algebraic models for structure spaces: Proofs • A. Homeomorphism groups of some stratified spaces • B. Controlled homeomorphism groups • C. $K$-theory of pairs and diagrams • D. Corrections and Elaborations Permission – for use of book, eBook, or Journal content Please select which format for which you are requesting permissions.
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Math Test #1 Solutions (a) (3 + 2x)(1 - 3x) = (6x - 1)(1 - x), (b) x(10x - 11) = 6 Solutions: (a) Expand both sides of the equation to obtain: (3+2x)(1-3x) = (6x-1)(1-x) ^2 = 6x-6x^2-1+x and so the solution written in set notation is (b) Using the distributive property again we have x(10x - 11) = 6 ^2 - 11x = 6 ^2 - 11x - 6 = 0: Since the factorization method is not easy to apply here and also the method of completing the square is still not very easy we apply the quadratic formula a = 10, b = -11, c = -6: (c) Eliminating the denominators we have then after checking we see the only solution is Solutions: (a) Squaring both sides, after isolating the radical ex- pression, we get just an implication but good enough to see how the solutions might look like: or x = 30 and after checking we see that (b) Using the same method we get 32x^2 - x - 3 = 0. Using the quadratic formula of which gives only one solution relative to the water in going 8 miles upstream and then returning. The total time for the trip was 1 hours. Use this information to find the speed of the current. (b) If instead of knowing that it took 1 hours for the round trip we know that it took 20 minutes more to go upstream than to go downstream what would be the speed of the current in this case? Solutions: (a) If we denote by x the speed of current in miles per hour we get that the time going upstream is time going downstream is hour round trip we obtain the equation: Solving this for x by adding to the same common denominator we which gives (b) Since 20 minutes is 1/3 of an hour. In this case the equation which implies This gives only the solution terval notation: Solutions: (a) The inequality is equivalent to which means (b) The inequality becomes (c) In this case we have similarly (d) Since we are dealing with positive numbers the inequality is equivalent to which gives x^2 + y^2 + 39x - 80y = 0. (b) What are the coordinates of the x-intercepts ? Solutions: (a) The equation can be written after completing the which gives the center (b) Making y = 0 in x^2+y^2+39x-80y = 0 we obtain x^2+39x = 0 which leads to two points of intersection: Similarly for the y-intercept we get: of this circle is included below: the point of coordinates (3,-2) and parallel to the line of equation 7x + 5y - 3 = 0. Solutions: The slope of the given line is m = equation of the line we are looking for is y - (-2) = (-7/5)(x - 3) Solutions: (a) We need to impose the condition: (b) In this case we have: 1. 2x - 1 ≥ 0 or to 2x - 1 = x^2 or 0 = x^2 - 2x + 1 (completing the square) (x - 1)^2 = 0. We have only one non-negative solution: x = 1.
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Translating and Solving Word Problems and Applications Learning Outcomes • Translate word phrases to equations and solve • Translate and solve applications In previous chapters, we translated word phrases into equations. This skill will help you when you solve word problems. Previously, you translated phrases into expressions, and now we will translate phrases into mathematical equations so we can solve them. But before we get started, let’s define some important terminology: • variables: Variables are symbols that stand for an unknown quantity, often represented with letters, like [latex]x, y[/latex], or [latex]z[/latex]. • coefficient: Sometimes a variable is multiplied by a number. This number is called the coefficient of the variable. For example, the coefficient of [latex]3x[/latex] is [latex]3[/latex]. • term: A single number, or variables and numbers connected by multiplication. For example, [latex]-4, 6x[/latex], and [latex]x^2[/latex] are all terms. • expression: Groups of terms connected by addition and subtraction. For example, [latex]2x^2-5[/latex] is an expression. • equation: An equation is a mathematical statement that two expressions are equal. An equation will always contain an equal sign with an expression on each side. Think of an equal sign as meaning “the same as”. Some examples of equations are [latex]y = mx +b[/latex], [latex]\frac{3}{4}r = v^{3} - r[/latex], and [latex]2(6-d) + f(3 +k) = \frac{1}{4}d[/latex] . The following figure shows how coefficients, variables, terms, and expressions all come together to make equations. In the equation [latex]2x-3^2=10x[/latex], the variable is [latex]x[/latex], a coefficient is [latex]10[/latex], a term is [latex]10x[/latex], and an expression is [latex]2x-3^2[/latex]. Translate Phrases into Equations The first step in translating phrases into equations is to look for the word (or words) that translate(s) to the equal sign. The table below reminds us of some of the words that translate to the equal sign. Equals (=) is is equal to is the same as the result is gives was will be Let’s review the steps we used to translate a sentence into an equation. Translate a word sentence to an algebraic equation. 1. Locate the “equals” word(s). Translate to an equal sign. 2. Translate the words to the left of the “equals” word(s) into an algebraic expression. 3. Translate the words to the right of the “equals” word(s) into an algebraic expression. In our first example, we will translate and solve a one-step equation. Translate and solve: five more than [latex]x[/latex] is equal to [latex]26[/latex]. Five more than [latex]x[/latex] [latex]\Rightarrow\quad{x+5}[/latex]is equal to [latex]\Rightarrow\quad{=}[/latex] Translate. [latex]26[/latex] [latex]\Rightarrow\quad{26}[/latex] Subtract 5 from both sides. [latex]x+5\color{red}{-5}=26\color{red}{-5}[/latex] Simplify. [latex]x=21[/latex] Check:Is [latex]26[/latex] five more than [latex]21[/latex] ? The solution checks. Translate and solve: The difference of [latex]5p[/latex] and [latex]4p[/latex] is [latex]23[/latex]. Watch this video for more examples of how to translate a phrase into an equation, then solve it. Translate and Solve Applications In most of the application problems we solved earlier, we were able to find the quantity we were looking for by simplifying an algebraic expression. Now we will be using equations to solve application problems. We’ll start by restating the problem in just one sentence, then we’ll assign a variable, and then we’ll translate the sentence into an equation to solve. When assigning a variable, choose a letter that reminds you of what you are looking for. The Robles family has two dogs, Buster and Chandler. Together, they weigh [latex]71[/latex] pounds. Chandler weighs [latex]28[/latex] pounds. How much does Buster weigh? Devise a problem-solving strategy. 1. Read the problem. Make sure you understand all the words and ideas. 2. Identify what you are looking for. 3. Name what you are looking for. Choose a variable to represent that quantity. 4. Translate into an equation. It may be helpful to restate the problem in one sentence with all the important information. Then, translate the English sentence into an algebra equation. 5. Solve the equation using good algebra techniques. 6. Check the answer in the problem and make sure it makes sense. 7. Answer the question with a complete sentence. Let’s take a look at the problem-solving strategy in action. Shayla paid $[latex]24,575[/latex] for her new car. This was $[latex]875[/latex] less than the sticker price. What was the sticker price of the car? Now you can try translating an equation from a statement that represents subtraction. In the following video, you will see another example of how to translate a phrase into an equation and solve.
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Post Correspondence Problem - (Theory of Recursive Functions) - Vocab, Definition, Explanations | Fiveable Post Correspondence Problem from class: Theory of Recursive Functions The Post Correspondence Problem (PCP) is a decision problem in theoretical computer science that involves finding a sequence of pairs of strings such that when the first components of these pairs are concatenated, they equal the concatenation of the second components. This problem is significant because it was one of the first problems proven to be undecidable, meaning there is no algorithm that can solve all instances of it. PCP is closely related to the concepts of formal languages and Turing machines, serving as a critical example in demonstrating the limits of computational solvability. congrats on reading the definition of Post Correspondence Problem. now let's actually learn it. 5 Must Know Facts For Your Next Test 1. The Post Correspondence Problem was introduced by Emil Post in 1946 and serves as a fundamental example in computability theory. 2. PCP is undecidable because there is no general algorithm that can solve all instances; some configurations will always lead to an infinite search without resolution. 3. The problem can be represented using pairs of strings, which can be visualized as tiles where each tile has a top and bottom string. 4. PCP has connections to other undecidable problems, such as the halting problem, illustrating broader implications in theoretical computer science. 5. A specific instance of PCP can often be constructed from finite automata or context-free grammars, demonstrating its relation to formal languages. Review Questions • How does the Post Correspondence Problem illustrate the concept of undecidability in computational theory? □ The Post Correspondence Problem exemplifies undecidability by showing that no algorithm can be constructed to solve every possible instance of the problem. This means that for some configurations, it is impossible to determine whether a solution exists, reflecting limits on what can be computed. As a result, PCP provides a concrete example used to understand the boundaries of algorithmic solvability in computer science. • Discuss the implications of the Post Correspondence Problem on understanding Turing machines and their limitations. □ The Post Correspondence Problem highlights significant limitations of Turing machines by showcasing a specific scenario where these machines cannot provide an answer. Since PCP is undecidable, it shows that even powerful computational models like Turing machines have boundaries beyond which they cannot operate. This relationship reinforces key concepts in computability theory about what problems are solvable and which are fundamentally beyond algorithmic reach. • Evaluate how the Post Correspondence Problem relates to other known undecidable problems like the halting problem and its broader impact on theoretical computer science. □ The Post Correspondence Problem is crucial in connecting with other undecidable problems such as the halting problem, as both illustrate fundamental limits within computation. By analyzing PCP alongside these problems, researchers gain insights into the nature of algorithmic computation and its boundaries. This understanding impacts theoretical computer science by establishing criteria for recognizing undecidable problems and shaping future research directions aimed at exploring new computational theories and models. "Post Correspondence Problem" also found in: © 2024 Fiveable Inc. All rights reserved. AP® and SAT® are trademarks registered by the College Board, which is not affiliated with, and does not endorse this website.
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12.8.1 Exercise: Projected Available Inventory Calculation Intended learning outcomes: Calculate projected available inventory by completing a grid. This exercise refers to Subsection 12.1.2. Complete the grid in Figure 12.8.1.1. Fig. 12.8.1.1 Projected available inventory calculation a. What is the available inventory without any restrictions along the time axis? b. What is the additional available inventory after order 102 9538? c. Which receipt could be deferred? d. Furthermore, the following orders are planned: • Customer order ID 104 2158 of 500 units on January 20 • Stock replenishment order ID 104 3231 of 500 units on January 22 Does this situation lead to a problem? If so, how can it be solved? a. 50 b. 300 (= 350 – 50) c. Stock replenishment order ID 101 2897 could be deferred to Jan. 14. d. Yes, there will not be enough available inventory on Jan. 20. Expedit­ing order ID 104 3231 by at least two days could solve this problem. Course section 12.8: Subsections and their intended learning outcomes
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Calculation of the flow rate of a rock-ramp pass The calculation of the flow rate of a rock-ramp pass corresponds to the implementation of the algorithm and the equations present in Cassan et al. (2016)^1. General calculation principle After Cassan et al., 2016^1 There are three possibilities: • the submerged case when \(h \ge 1.1 \times k\) • the emergent case when \(h \le k\) • the quasi-emergent case when \(k < h < 1.1 \times k\) In the quasi-emergent case, the calculation of the flow corresponds to a transition between emergent and submerged case formulas: \[Q = a \times Q_{submerge} + (1 - a) \times Q_{emergent}\] with \(a = \dfrac{h / k - 1}{1.1 - 1}\) Submerged case The calculation is done by integrating the velocity profile in and above the macro-roughnesses. The calculated velocities are the temporal and spatial averages per plane parallel to the bottom. In macro-roughnesses, velocities are obtained by double averaging the Navier-Stokes equations in uniform regime with a mixing length model for turbulence. Above the macro-roughnesses, the classical turbulent boundary layer analysis is maintained. The velocity profile is continuous at the top of the macro-roughnesses and is dependent on the boundary conditions set by the hydraulics: • velocity at the bottom (without turbulence) in m/s: \[u_0 = \sqrt{2 g S D (1 - \sigma C)/(C_d C)}\] • total shear stress at the top of the roughnesses in m/s: \[u_* = \sqrt{gS(h-k)}\] The average bed velocity is given by integrating the flows between and above the blocks: \[\bar{u} = \frac{Q_{inf} + Q_{sup}}{h}\] with respectively \(Q_{inf}\) and \(Q_{sup}\) the unit flows for the part in the canopy and the part above the canopy. Calculation of the unit flow rate Q[inf] in the canopy The flow in the canopy is obtained by integrating the velocity profile (Eq. 9, Cassan et al., 2016): \[Q_{inf} = \int_{0}^1 u(\tilde{z}) d \tilde{z}\] \[u(\tilde{z}) = u_0 \sqrt{\beta \left( \frac{h}{k} -1 \right) \frac{\sinh(\beta \tilde{z})}{\cosh(\beta)} + 1}\] \[\beta = \sqrt{(k / \alpha_t)(C_d C k / D)/(1 - \sigma C)}\] \[C_d = C_{x} f_{h_*}(h_*)\] and \(\alpha_t\) obtained by solving the following equation: \[\alpha_t u(1) - l_0 u_* = 0\] \[l_0 = \min \left( s, 0.15 k \right)\] \[s = D \left( \frac{1}{\sqrt{C}} - 1 \right)\] Calculation of the unit flow Q[sup] above the canopy \[Q_{sup} = \int_k^h u(z) dz\] with (Eq. 12, Cassan et al., 2016) \[u(z) = \frac{u_*}{\kappa} \ln \left( \frac{z - d}{z_0} \right)\] with (Eq. 14, Cassan et al., 2016) \[z_0 = (k - d) \exp \left( {\frac{-\kappa u_k}{u_*}} \right)\] and (Eq. 13, Cassan et al., 2016) \[ d = k - \frac{\alpha_t u_k}{\kappa u_*}\] which gives \[Q_{sup} = \frac{u_*}{\kappa} \left( (h - d) \left( \ln \left( \frac{h-d}{z_0} \right) - 1\right) - \left( (k - d) \left( \ln \left( \frac{k-d}{z_0} \right) - 1 \right) \right) \right)\] Emerging case The calculation of the flow rate is done by successive iterations which consist in finding the flow rate value allowing to obtain the equality between the flow velocity \(V\) and the average velocity of the bed given by the equilibrium of the friction forces (bottom + drag) with gravity: \[u_0 = \sqrt{\frac{2 g S D (1 - \sigma C)}{C_d f_F(F) C (1 + N)}}\] \[N = \frac{\alpha C_f}{C_d f_F(F) C h_*}\] \[\alpha = 1 - (a_y / a_x \times C)\] Formulas used Bulk velocity V \[V = \frac{Q}{B \times h}\] Average speed between blocks V[g] From Eq. 1 Cassan et al (2016)^1 and Eq. 1 Cassan et al (2014)^2: \[V_g = \frac{V}{1 - \sqrt{(a_x/a_y)C}}\] Drag coefficient of a single block C[d0] \(C_{d0}\) is the drag coefficient of a block considering a single block infinitely high with \(F << 1\) (Cassan et al, 2014^2). Block shape Cylinder "Rounded face" shape Square-based parallelepiped "Flat face" shape Value of \(C_{d0}\) 1.0 1.2-1.3 2.0 2.2 When establishing the statistical formulae for the 2006 technical guide (Larinier et al. 2006^4), the definition of the block shapes to be tested was based on the use of quarry blocks with neither completely round nor completely square faces. The so-called "rounded face" shape was thus not completely cylindrical, but had a trapezoidal bottom face (seen in plan). Similarly, the "flat face" shape was not square in cross-section, but also had a trapezoidal bottom face. These differences in shape between the "rounded face" and a true cylinder on the one hand, and the "flat face" and a true parallelepiped with a square base on the other hand, result in slight differences between them in the shape coefficients \(C_{d0}\). Block shape coefficient σ Cassan et al. (2014)^2, et Cassan et al. (2016)^1 define \(\sigma\) as the ratio between the block area in the \(x,y\) plane and \(D^2\). For the cylindrical form of the blocks, \(\sigma\) is equal to \(\pi / 4\) and for a square block, \(\sigma = 1\). Ratio between the average speed downstream of a block and the maximum speed r The values of \(r\) depends on the block shapes (Cassan et al., 2014^2 et Tran et al. 2016 ^3): • round : \(r_Q=1.1\) • "rounded face" shape : \(r=1.2\) • square-based parallelepiped : \(r=1.5\) • "flat face" shape : \(r=1.6\) Cassiopée implements a formula depending on \(C{d0}\): \[ r = 0.4 C_{d0} + 0.7 \] Froude F \[F = \frac{V_g}{\sqrt{gh}}\] Froude-related drag coefficient correction function f[F](F) If \(F < 1\) (Eq. 19, Cassan et al., 2014^2): \[f_F(F) = \min \left( \frac{r}{1- \frac{F_{g}^{2}}{4}}, \frac{1}{F^{\frac{2}{3}}} \right)^2\] otherwise \(f_F(F) = 1\) because a torrential flow upstream of the blocks is theoretically impossible because of the hydraulic jump caused by the downstream block. Maximum speed u[max] According to equation 19 of Cassan et al, 2014^2 : \[ u_{max} = V_g \sqrt{f_F(F)} \] Drag coefficient correction function linked to relative depth f[h*](h[*]) The equation used in Cassiopeia differs slightly from equation 20 of Cassan et al. 2014^2 and equation 6 of Cassan et al. 2016^1. This formula is a fit to the experimental measurements on circular blocks used in Cassan et al. 2016^1: \[ f_{h_*}(h_*) = (1 + 1 / h_*^{2}) \] Coefficient of friction of the bed Cf If \(k_s < 10^{-6} \mathrm{m}\) then we use Blasius' formula \[C_f = 0.3164 / 4 * Re^{-0.25}\] \[Re = u_0 \times h / \nu\] Else (Eq. 3, Cassan et al., 2016 d'après Rice et al., 1998^5) \[C_f = \frac{2}{(5.1 \mathrm{log} (h/k_s)+6)^2}\] • \(\alpha\): ratio of the area affected by the bed friction to \(a_x \times a_y\) • \(\alpha_t\): length scale of turbulence in the block layer (m) • \(\beta\): ratio between drag stress and turbulence stress • \(\kappa\): Von Karman constant = 0.41 • \(\sigma\): ratio between the block area in the plane X,y et \(D^2\) • \(a_x\): cell width (perpendicular to the flow) (m) • \(a_y\): length of a cell (parallel to the flow) (m) • \(B\): pass width (m) • \(C\): blocks concentration • \(C_d\): drag coefficient of a block under current flow conditions • \(C_{d0}\): drag coefficient of a block considering an infinitely high block with \(F \ll 1\) • \(C_f)\): bed friction coefficient • \(d\): displacement in the zero plane of the logarithmic profile (m) • \(D\): width of the block facing the flow (m) • \(F\): Froude number based on \(h\) and \(V_g\) • \(g\): acceleration of gravity = 9.81 m.s^-2 • \(h\): average depth (m) • \(h_*\): dimensionless depth (\(h / D\)) • \(k\): useful block height (m) • \(k_s\): roughness height (m) • \(l_0\): length scale of turbulence at the top of the blocks (m) • \(N\): ratio between bed friction and drag force • \(Q\): flow (m^3/s) • \(S\): pass slope (m/m) • \(u_0\): average bed speed (m/s) • \(u_*\): shear velocity (m/s) • \(V\): flow velocity (m/s) • \(V_g\): velocity between blocks (m/s) • \(s\): minimum distance between blocks (m) • \(z\): vertical position (m) • \(z_0\): hydraulic roughness (m) • \(\tilde{z}\): dimensionless stand \(\tilde{z} = z / k\)
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Radiation Modeling with Participating Media In AcuSolve enclosure or surface to surface radiation models, effect of media between surfaces is ignored. This assumption is acceptable when you are dealing with lower temperature fluids. Nonetheless, while you have semi-transparent media like glass or high temperature gases like in flames, effect of media should be considered in heat transfer analysis. Two models are available in AcuSolve: A simple one equation P[1] model and more detailed but expensive discrete ordinates (DO) model. Radiative Transfer Equation (RTE) Radiative energy balance in a participating media is governed by the following integro-differential equation, known as the radiative transfer equation (RTE): $\underset{\text{rate}\text{}\text{}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{of}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{change}}{\underbrace{\Omega ·abla I\left(r,\Omega \ right)}}+\underset{\text{Absorption}}{\underbrace{\kappa I\left(r,\Omega \right)}}+\underset{\text{Scattering loss}}{\underbrace{{\sigma }_{s}I\left(r,\Omega \right)}}=\underset{\text{Emission}}{\ underbrace{f\left(r\right)}}+\underset{\text{Scattering addition}}{\underbrace{\frac{{\sigma }_{s}}{4\pi }\underset{4\pi }{\int }I\left(r,{\Omega }^{\prime }\right)\varphi \left({\Omega }^{\prime },\ Omega \right)d{\Omega }^{\prime }}}$ is the radiation intensity, is the spatial vector, is the unit directional vector, ${\Omega }^{\prime }$ is the scattering directional vector, is the absorption coefficient, and ${\sigma }_{s}$ is the scattering coefficient, is an emission, the refractive index, and is the Stefan-Boltzman constant (5.67 × 10 W m ). T is the temperature (K). • P1 radiation model • Discrete Ordinates (DO) model P1 Radiation Model The P[1] model is the lowest order P[N] (spherical harmonics) type radiation model. The method reduces the five independent variables of the Radiative Transfer Equation (RTE) into a PDE that is relatively simple in comparison. The model is the most computationally efficient of the radiation models in AcuSolve, but it can lose accuracy, under certain conditions, for optically thin media. It performs best in scenarios where the radiative intensity is near isotropic. Governing Equation P[1] approximation and assumptions The P[1] model is derived from the general P[N] formulation (a spherical harmonic series expansion of the radiative intensity for the angular variable) by assuming that the series is limited to four terms and integrating over all solid angles. From the first harmonic in the series approximation, the divergence of the radiative flux ( $q$) can be derived by integrating the RTE over all solid angles as $abla \cdot q=\kappa \left(4{n}^{2}\sigma {T}^{4}-G\right)$ where $G$ is the incident radiation $\frac{{\sigma }_{s}}{4\pi }\underset{4\pi }{\int }I\left(r,{\Omega }^{\prime }\right)\varphi \left({\Omega }^{\prime },\Omega \right)d{\Omega }^{\prime }$, $\ kappa$ is the absorption coefficient, $n$ the refractive index, and $\sigma$ is the Stefan-Boltzman constant (5.67 × 10^-8 W m^-2 K^-4), $q$ is the radiative flux. Additionally, a second vector equation can be derived from the other three harmonic terms for the radiative flux $q=\Gamma \text{}\text{\hspace{0.17em}}abla G$ where $\Gamma$ is the diffusion coefficient. By taking the divergence of (2), and substituting into this the right hand side of equation (1), leads to elimination of the heat flux. The final diffusion reaction equation describing the transport of incident radiation is given by $\Gamma \text{\hspace{0.17em}}{abla }^{2}G-\kappa \left(4{n}^{2}\sigma {T}^{4}-G\right)=0$ where $\Gamma$ is a diffusion coefficient given by $\Gamma =\frac{1}{3\left(\kappa +{\sigma }_{s}\right)-{A}_{1}{\sigma }_{s}}$ where $\kappa$ is the absorption coefficient, ${\sigma }_{s}$ is the scattering coefficient, and ${A}_{1}$ is the linear-anisotropic phase function. Anisotropic scattering The implementation in AcuSolve includes the ability to model linear anisotropic scattering $\varphi \left({\Omega }^{\prime }\cdot \Omega \right)=1+{A}_{1}{\Omega }^{\prime }\cdot \Omega$ is the unit directional vector in the scattering direction, ${\Omega }^{\prime }$ is the unit directional vector in the incident radiation direction. This term is included in the scattering term of the RTE along with the four term spherical harmonic expansion of the radiative intensity field (the first term being isotropic and the other three anisotropic). The final form of the P approximation in equation (3) includes this term. The values of the phase function A have the following meaning: • A[1] = 1: More radiation is scattered in the forward direction • A[1]= -1: More radiation is scattered in the backward direction • A[1]= 0: Isotropic scattering Coupling to energy equation The source term in equation (1) can be substituted into the energy equation as a negative source since a local increase in radiative heat flux is due to a local decrease in thermal energy. the default stagger sequence when the P radiative heat transfer solver is enabled is: 1. Solve the energy equation with source term ( κ(4n^2σT^4-G) ), where G is zero for the time step 2. Pass temperature to radiation solver and solve equation (3) 3. Repeat until converged Boundary Conditions Marshak's boundary condition Based on the assumption that the walls are diffused gray surfaces, that is, independent of wavelength, the appropriate wall boundary condition is $n\cdot abla G=\frac{\epsilon }{2\left(2-\epsilon \right)}\left(4{n}^{2}\sigma {T}_{w}^{4}-{G}_{w}\right)$ where $\epsilon$ is the surface emissivity, ${T}_{w}^{}$ is the wall temperature, and ${G}_{w}$ the wall incident radiation. The boundary radiative heat flux can be calculated from the incident radiation and the temperature at the wall: ${q}_{w}=\frac{\epsilon }{2\left(2-\epsilon \right)}\left({G}_{w}-4{n}^{2}\sigma {T}_{w}^{4}\right)$ Discrete Ordinates (DO) Model Governing equation The governing equation is the radiative transfer equation limited to a finite number of directions (or ordinates) $\Omega ·abla I\left(r,\Omega \right)+\kappa I\left(r,\Omega \right)+{\sigma }_{s}I\left(r,\Omega \right)=f\left(r\right)+\frac{{\sigma }_{s}}{4\pi }\underset{4\pi }{\int }I\left(r,{\Omega }^{\prime }\right)\varphi \left({\Omega }^{\prime },\Omega \right)d{\Omega }^{\prime }$ $f\left(r\right)=\kappa {I}_{b}\left(r\right)=\kappa \frac{\sigma {T}^{4}}{\pi }$ Scattering term (source term) The integral over all the directions in equation (1) is replaced by a numerical quadrature for $N$ different ordinate directions ( ${\Omega }_{i}$) $S=\frac{{\sigma }_{s}}{4\pi }\sum _{j=1}^{N}{w}_{j}I\left(r,{\Omega }_{j}\right)\varphi \left({\Omega }_{j},{\Omega }_{i}\right)$ The phase function, $\varphi \left({\Omega }_{j},{\Omega }_{i}\right)$, is given by $\varphi \left({\Omega }^{\prime }\cdot \Omega \right)=1+{A}_{1}\left({\Omega }^{\prime }\cdot \Omega \right)$ Boundary conditions (RTE) Diffused surface If a surface emits and reflects diffusely, the exiting intensity is directionally independent and is given by $I\left({r}_{w},{\Omega }_{i}\right)=\epsilon \left({r}_{w}\right){I}_{b}\left({r}_{w}\right)+\frac{1-\epsilon \left({r}_{w}\right)}{\pi }\sum _{n\cdot {\Omega }_{j}<0}{w}_{j}I\left({r}_{w},{\Omega } _{j}\right)\left|n\cdot {\Omega }_{j}\right|,\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}} \text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}n\cdot {\Omega }_{j}>0\text{\hspace{0.17em}}$ Specular and diffuse surface The diffused fraction defines the proportion of reflected radiation intensity at a surface which is diffused, that is, the reflection may also have a specular component. If the radiation intensity reflection coefficient at the surface is defined by $\rho ={\rho }^{S}+{\rho }^{D}=1-\epsilon$ then the diffused reflection coefficient, ${\rho }^{D}$, is defined in terms of the diffused fraction $\alpha$ and the emissivity of the surface by ${\rho }^{D}=\alpha \left(1-\epsilon \right)$ and the specular reflection coefficient ${\rho }^{S}=\left(1-\alpha \right)\left(1-\epsilon \right)$ If $\alpha$=1, then the reflection at the surface is completely diffused. If $\alpha$=0 then the reflection is specular. The outgoing intensity, I, at the surface in terms of the above two reflection coefficients is given by $I\left({r}_{w},{\Omega }_{i}\right)=\epsilon \left({r}_{w}\right){I}_{b}\left({r}_{w}\right)+\frac{{\rho }^{D}\left({r}_{w}\right)}{\pi }\sum _{n\cdot {\Omega }_{j}<0}^{N}{w}_{j}I\left({r}_{w},{\ Omega }_{j}\right)\left|n\cdot {\Omega }_{j}\right|+{\rho }^{S}\left({r}_{w}\right)I\left({r}_{w},{\Omega }_{j}\right),\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace {0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}n\cdot {\Omega }_{j}>0\text{\hspace where the first terms represent emission from the surface, the second term the diffused component incoming radiation heat flux and the third the specular component. The diffused component represents a sum over all radiation intensities along ordinates that are incident to the surface (that is, a hemisphere of incoming radiation to the surface); $n$ is the normal into the domain and ${\Omega }_ {j}$ the jth ordinate direction. The ordinate direction ( $\Omega$), the total number of ordinate directions ( $N$) and the weights ( $w$) are automatically defined by the order of the radiation_quadrature (S2, S4, S6, S8 & S10). The specular ordinate direction ( ${\Omega }_{S}$) is the direction that the radiation intensity must strike the surface to reflect in a specular fashion along the outgoing ordinate direction, ${\Omega }_{i}$, and is given by ${\Omega }_{s}={\Omega }_{i}-2\left({\Omega }_{i}\cdot n\right)n$ which means the angle that incident radiation intensity strikes the surface equals the angle of reflection. Boundary conditions (Energy equation) Interface and outflow/inflow boundary conditions At an opaque interface between participating and non-participating media or outflows/inflows a radiative heat flux must be added to the boundaries in the energy equation. This flux is given by ${Q}_{rad}=\epsilon \left(\sum _{{\Omega }_{j}\cdot n>0}{w}_{j}I\left({r}_{w}\right)\left|{\Omega }_{j}\cdot n\right|-{n}^{2}\sigma {T}_{w}^{4}\right)\text{\hspace{0.17em}}$ For an opening, that is, outflow or inflow, the black body intensity used in the calculation of the outgoing intensity at the surface is based on the opening temperature of the surrounding: ${I}_{b}=\frac{\kappa {n}^{2}\sigma {T}_{open}^{4}}{\pi }$ while for an interface it is based on the current temperature solution. Output metrics Two directionally integrated output metrics can be derived from the radiative intensities: incident radiation and radiative heat flux. Incident radiation Incident radiation is the total intensity impinging on a point from all directions and is given by $G=\sum _{i=1}^{N}{w}_{i}I\left(r\right)$ where $I$ is the intensity in direction i, $N$ the number of ordinates, $w$ the weights. Interface Between two Semitransparent Media At the interface between two semitransparent media (referred to as medium 1 and medium 2 below), radiative intensity is both transmitted and reflected. The proportion of transmitted and reflected intensity at the interface depends on: the refractive indices (n[1], n[2]) of the two media; the incident angle that radiative intensity strikes the surface; and the diffuse fraction of the surface. Reflection and transmission for specular interfaces Reflection at an interface is governed by the angle of incidence of a radiative intensity and the refractive indices of the two media. The image below shows the different refracted and reflected rays between two media. Figure 1. Reflected and refracted directions at the interface between two participating media of different refractive indices (n[1] < n[2]). For the medium of higher refractive index if the incoming direction is greater than the critical angle total internal reflection occurs (no transmission occurs across the interface). This is represented by the gray dashed lines in the image. The cosine of the incident angle for the incoming ordinate is given by $\mathrm{cos}{\theta }_{1}={\Omega }_{I}^{1}\cdot n$ where $n$ is the outward facing normal direction at the interface (towards the second medium) and ${\Omega }_{I}^{1}{}_{1}$ is the unit direction of incoming radiative intensity to the surface, given ${\Omega }_{I}^{1}{}_{1}={\Omega }_{R}^{1}-2\left({\Omega }_{R}^{1}\cdot n\right)n$ where ${\Omega }_{I}^{1}{}_{1}$ is the unit reflected ordinate direction vector and also represents the current ordinate direction being solved. The equivalent calculation can also be performed for medium two. Radiative intensity that is transmitted into a second medium undergoes refraction governed by Snell's law, ${n}_{1}\mathrm{sin}{\theta }_{1}={n}_{2}\mathrm{sin}{\theta }_{2}$ where n[1] and n[2] are the refractive indices of mediums. θ[1] and θ[2] are the angles of incidence and refraction of radiative intensity relative to the interface normal, respectively. This can also be represented in vector form by ${n}_{1}\left(n×{\Omega }_{I}^{1}\right)={n}_{2}\left(n×{\Omega }_{R}^{2}\right)$ The incoming direction vector in medium two for a ray refracted from medium two to one is given by ${\Omega }_{I}^{2}=\frac{{n}_{1}}{{n}_{2}}{\Omega }_{R}^{1}+\left(\frac{{n}_{1}}{{n}_{2}}\mathrm{cos}{\theta }_{1}-\sqrt{1-{\left(\frac{{n}_{1}}{{n}_{2}}\right)}^{2}\left(1-\mathrm{cos}{\theta }_{1}\ The above expression is valid providing the expression under the radicand is greater than zero; otherwise total internal reflection occurs. The actual reflected and refracted directions differ slightly from the calculated direction since these directions will unlikely coincide with a discrete ordinate direction. Since the number of directions is governed by the order of radiation quadrature, higher quadrature orders are more accurate for interface problems. Figure 2. Octant of angular quadrature with transmission direction. The calculated transmission direction (Ω[T]) is shifted to match the nearest quadrature point. Depending on the refractive indices of the two media and the angle of incidence, θ[1], the proportion of radiation intensity that is reflected or transmitted will vary. If ${n}_{1}<{n}_{2}$, then the radiative intensity in medium one will be partially reflected and partially transmitted into a cone defined by the critical angle, θ[c], which is given by: ${\theta }_{c}={\mathrm{sin}}^{-1}\frac{{n}_{1}}{{n}_{2}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\ and defines a cone in 3D (see the images below). The critical angle is defined by a ray that grazes the surface on the side of medium 2 and is transmitted exactly at the critical angle in medium 1. As an example, if ${n}_{1}>{n}_{2}$ and ${\theta }_{1}>{\theta }_{c}$, the direction of incoming radiation is greater than the critical angle and total internal reflection will occur. This means the incoming ray is reflected at the same angle of incidence and no transmission occurs. Figure 3. Total internal reflection of intensity rays at an interface ( ${n}_{1}>{n}_{2}$). The critical angle cone defines the angle outside which total internal reflection occurs, that is, ${\theta }_{1}>{\theta }_{c}$. However, if ${\theta }_{1}<{\theta }_{c}$, the outgoing intensity will include transmission from medium 2 to medium 1, as shown in the image below. Figure 4. Reflection and transmission of intensity rays at an interface ( ${n}_{1}>{n}_{2}$) for ${\theta }_{1}<{\theta }_{c}$ At the interface, the intensity is partially reflected and transmitted into the other medium if ${\theta }_{1}<{\theta }_{c}$ or the outgoing direction is in medium 2. The reflected proportion from $ {\Omega }_{I}^{1}\to {\Omega }_{R}^{1}$, or reflectance, is given by ${r}_{12}=\frac{1}{2}{\left(\frac{{n}_{1}\mathrm{cos}{\theta }_{1}-{n}_{2}\mathrm{cos}{\theta }_{2}}{{n}_{1}\mathrm{cos}{\theta }_{1}+{n}_{2}\mathrm{cos}{\theta }_{2}}\right)}^{2}+\frac{1}{2}{\left(\ frac{{n}_{2}\mathrm{cos}{\theta }_{1}-{n}_{1}\mathrm{cos}{\theta }_{2}}{{n}_{2}\mathrm{cos}{\theta }_{1}+{n}_{1}\mathrm{cos}{\theta }_{2}}\right)}^{2}$ and the transmitted proportion from the ${\Omega }_{I}^{2}\to {\Omega }_{R}^{1}$, or transmittance, is given by ${\tau }_{21}=1-{r}_{12}$ In the second medium, for the current scenario where ${n}_{2}>{n}_{1}$, if ${\theta }_{2}<{\theta }_{c}$ the radiative intensity is, as for medium one, partially reflected and partially transmitted. The reflection coefficient is as described above since ${r}_{21}={r}_{12}$. If ${\theta }_{2}>{\theta }_{c}$, then total internal reflection occurs and ${r}_{21}=1.0$ and ${\tau }_{12}=0.0$, meaning no transmission of radiative intensity into the second medium or from the first medium. This is shown in the image above with the gray dashed lines. From the above, the outgoing radiative intensity on the side one of the interface for the current ordinate direction, ${\Omega }_{R}^{1}$ is given by $I\left({\Omega }_{R}^{1}\right)={r}_{12}I\left({\Omega }_{I}^{1}\right)+{\tau }_{21}{\left(\frac{{n}_{1}}{{n}_{2}}\right)}^{2}I\left({\Omega }_{I}^{2}\right)$ where for medium one, the first term on the right-hand side represents the reflected intensity in medium one and the second term represents the transmitted intensity from medium two to one. For medium two, if the current ordinate direction is ${\Omega }_{R}^{2}$ then the intensity outgoing radiative intensity is given by $I\left({\Omega }_{R}^{2}\right)={r}_{21}I\left({\Omega }_{I}^{1}\right)+{\tau }_{12}{\left(\frac{{n}_{2}}{{n}_{1}}\right)}^{2}I\left({\Omega }_{I}^{1}\right)$ For ${n}_{2}<{n}_{1}$, the subscripts of the above analysis must be exchanged, and total internal reflection will now occur in medium one. Reflection and transmission for diffuse interfaces If the interface is diffused, for example, diffused_fraction = 1.0, the reflectivity of the interface is given by the hemispherically averaged reflectance: is the ratio of refractive indices. Note: ${n}_{1}$ always represents the medium with higher refractive index and ${n}_{2}$ the medium of lower refractive index. The transmission from medium one to two is given by ${\tau }_{D,12}=1-{r}_{D,12}$ For the reverse direction the reflectance and transmittance are given by: ${\tau }_{D,21}=1-\frac{1}{{n}^{2}}\left(1-{r}_{D,12}\right)$ ${\tau }_{D,21}=\frac{1}{{n}^{2}}{\tau }_{D,12}$ The incoming radiative intensity to the interface is given by the hemispherically averaged intensity for medium one and two: ${Q}_{1}=\sum _{j=1}^{N}{w}_{j}{I}_{j}\left|n\cdot {\Omega }_{j}\right|\left(n\cdot {\Omega }_{j}\right)\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\ hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{when}\text{\hspace{0.17em}}\text{\hspace{0.17em}}n\cdot {\Omega }_{j}>0$ ${Q}_{2}=\sum _{j=1}^{N}{w}_{j}{I}_{j}\left|n\cdot {\Omega }_{j}\right|\left(n\cdot {\Omega }_{j}\right)\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{\ hspace{0.17em}}\text{\hspace{0.17em}}\text{\hspace{0.17em}}\text{when}\text{\hspace{0.17em}}\text{\hspace{0.17em}}n\cdot {\Omega }_{j}\le 0$ where $n$ is the outward facing normal. From these fluxes, the outgoing radiative intensity at the wall for the current ordinate direction, $\Omega$, is given by $I\left(\Omega \right)={r}_{D,12}\frac{{Q}_{1}}{\pi }+{\tau }_{D,21}\frac{{Q}_{2}}{\pi }$ Reflection and Transmission for partially specular and partially diffuse interfaces For partially specular and partially diffuse interfaces 0.0 < diffused fraction < 1.0. Interfaces between semi-transparent media are typically not 100 percent diffused or specular and the diffuse fraction lies somewhere between zero and one. In this range the outgoing radiative intensity is treated as a linear combination of the specular and diffuse components, for example: $I\left(\Omega \right)=\left(1-\alpha \right){I}^{S}\left(\Omega \right)+\alpha {I}^{D}\left(\Omega \right)$ where $\alpha$ is the diffuse fraction, ${I}^{S}\left(\Omega \right)$ is the outgoing specular component of radiative intensity, and ${I}^{D}\left(\Omega \right)$ is the outgoing diffuse component of radiative intensity. For example, in medium one in the image above the components would be: ${I}^{S}\left(\Omega \right)=I\left({\Omega }_{R}^{1}\right)={r}_{12}I\left({\Omega }_{I}^{1}\right)+{\tau }_{21}{\left(\frac{{n}_{1}}{{n}_{2}}\right)}^{2}I\left({\Omega }_{I}^{2}\right)$ ${I}^{D}\left(\Omega \right)={r}_{D,12}\frac{{Q}_{1}}{\pi }+{\tau }_{D,21}\frac{{Q}_{2}}{\pi }$ Specular and diffuse interfaces For the interface equations to be applied weakly, I[w] must be applied in both mediums. If the current ordinate direction, $\Omega$, is outgoing from the interface in medium 1 then I[w] is equal to the proportion of radiative intensity reflected and transmitted. From the analysis in the previous section, this would be: ${I}_{w}={r}_{12}I\left({\Omega }_{I}^{1}\right)+{\tau }_{21}{\left(\frac{{n}_{1}}{{n}_{2}}\right)}^{2}I\left({\Omega }_{I}^{2}\right)$ In medium 2 since the current ordinate direction, $\Omega$, is incoming to the surface no boundary flux is added to the equation, that is, $\eta {\left({I}_{w},v\right)}_{\Gamma }=0$. Reflection and Transmission for diffuse interfaces of Type External Exchange of radiative intensity occurs for external surfaces when the medium inside the computational domain is semitransparent. That is the medium surrounding, which is not modeled using a computational mesh, participates in radiative transfer. For this case a mathematical model of external radiation is used. The model assumes that the surrounding medium has uniform radiative intensity in all directions, that is, the radiative flux is isotropic. The isotropic radiative intensity is given by the following blackbody source: ${I}_{\text{ext}}=\frac{{\epsilon }_{\text{ext}}\sigma {T}_{\text{ext}}^{4}}{\pi }$ where ${\epsilon }_{\text{ext}}$ is the external emissivity and is set to one, $\sigma$ is the Stefan-Boltzmann constant, ${T}_{\text{ext}}$ is the temperature of the surrounding medium. At the external interface ${I}_{\text{ext}}$ is transferred into the medium. This condition can only be applied to boundaries as the interface is only modeled mathematically.
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factors affecting size of product from ball mill WEBAug 19, 2022 · In all instances, Cyclone performance can be affected by six factors, including: Size. Flow rate. Inlet area. Vortex finder diameter. Underflow diameter. Length of Cyclone. Let's take a look at how and why these factors affect Cyclone performance. 1. WhatsApp: +86 18838072829 WEBJan 1, 2022 · The size of grinding media is the primary factor that affects the overall milling efficiency of a ball mill ( power consumption and particle size breakage). This article tackles the lack of a ... WhatsApp: +86 18838072829 WEBJan 12, 2021 · The hammer mill works by impact between rapidly moving hammers and powder materials, while the ball mill uses impact and attrition between rapidly moving balls inside a rotating cylinder. Both provide descriptions of their basic principles, constructions, operating parameters, advantages and disadvantages. Read more. 1 of 53. Download . WhatsApp: +86 18838072829 WEBApr 22, 2019 · 5. 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Answer: Colloid mill WhatsApp: +86 18838072829 WEBDec 5, 2019 · A. Ashwini Shewale. Size reduction, or comminution, is the process of reducing substances into smaller pieces or fine powder. It increases the surface area, allowing for faster dissolution and extraction of active constituents. The degree of size reduction depends on factors like the drug's hardness, solvent used, and extraction . WhatsApp: +86 18838072829 WEBMay 1, 2020 · The production capacity of the largescale ball mill in the concentrator is a crucial factor affecting the subsequent separation and the economic benefits of the operation. WhatsApp: +86 18838072829 WEBFeb 1, 2018 · It was found that the ball mill consumed kWh/t energy to reduce the F 80 feed size of lm to P 80 product size of lm while stirred mill consumed kWh/t of energy to produce ... WhatsApp: +86 18838072829
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Longest Common Substring - Scaler Blog Problem statement: Given two strings string_1 and string_2 of size N and M respectively. Find the length of the longest common substring. The problem statement says that we are given two strings and we need to return the length of the longest common substring from the given strings. For example: string_1 = "fishes" string_2 = "fisherman" The longest common substring from string_1 and string_2 of lengths N = 6 and M = 9 respectively, is “fishe” and it is of length 5. Therefore the output is 5. Let’s discuss various ways for obtaining the length of the longest common substring. • 1<N<10^4 • 1<M<10^4 • where N and M are the lengths of the string_1 and string_2 respectively. What is a Substring? A substring is the contiguous sequence of characters within the given string. For example, “string” is a substring of “substring” where the order remains the same. For instance, Consider given input string is “abce” The substrings of the given string are: [a, b, c, e, ab, bc, ce, abc, bce, abce] A substring is contiguous, unlike a subset. Even the whole string is also considered to be one of the substrings of the given string. If N is the length of the given string, then the number of non-empty substrings from the given string is calculated by: Number of non-empty substrings = (N*(N+1))/2 Now as we have a basic understanding of substring let’s discuss the longest common substring. When two strings are given we need to find the length of the longest substring among these two strings. For instance, consider string_1 = "acevd" string_2 = "acenbvd" The common substrings among these two strings are: But the longest common substring is “ace” of length 3. Approach 1 (Brute Force Approach): • Generate all the substrings of string_1 and string_2 using three nested loops and store them in any data structure like an array. • For generating substrings we use three nested loops: □ The outer loop keeps track of starting character □ The second loop tracks all the characters on the right of the starting character as the ending character of the substring. □ The innermost/third loop appends the characters from the current starting index to the ending index. □ Therefore generating substrings takes a time complexity of O(N^2) where N is the length of the strings given. • Now from the substrings from the arrays, compare each substring of each string given and always carry forward the maximum length of the common substring formed from the given two strings. • Return the length of the longest common substring obtained. Since this is a naive approach we can just follow the algorithm and create all the substrings and compare for the length of the longest common substring from the given strings. Complexity Analysis: Time Complexity: O(N^2 * M^2) where N and M are the lengths of the string_1 and string_2 respectively. • For generating the substrings for string_1 and string_2 and comparing the formed substring one by one. • Here, string_1 has N^2 number of substrings and string_2 has M^2 number of strings. • For comparing all these formed substrings of string_1 and string_2, it takes O(N^2^ * M^2^). • Therefore the overall time complexity is O(N^2 * M^2) Space Complexity: O(N^2) + O(M^2) where N and M are the lengths of the string_1 and string_2 respectively. • For storing the substrings formed by each string. • string_1 takes O(N^2^) space and string_2 takes O(M^2^) and the sum of it is the space complexity. Approach 2: Efficient Approach using Recursion The above approach generates all the substrings of given strings and compares iteratively. Instead of redundantly creating all the substrings, we can just use recursion for comparing the strings. If the character at which we are checking is matched then, we increment the LCSCount by 1 else, making it 0. This algorithm uses recursion for finding the length of the longest common substring from the given two strings. The algorithm is as follows: • Start traversing from the end of the two strings with i and j for string_1 and string_2 respectively using recursion. • If the last characters of both the strings match then decrement both the iterators by 1 and increment LCSCount by 1 and store it in LCSCount_1. • Else, we need to consider the following two cases: □ Just decrement the iterator of string_1 and call the recursive function by updating LCSCount to 0 (because we missed the contiguous character of the string) and store it in LCSCount_2. □ The other way can be decrementing the iterator of string_2 by 1 and again calling the recursive function by updating the LCSCount to 0 and store it in LCSCount_3 • The base case of this function would be till iterator_1 and iterator_2 should be greater than zero, else return the LCSCount because once any one of the iterators becomes zero it means that we have come to an end of the string given and there can be no further longest common substring which will be formed so, we need to return the LCSCount. • The answer in this approach is the maximum of LCSCount1, LCSCount_2, LCSCount_3. This recursive relations can be better explained using: Below is the implementation of the above approach: C++ implementation for finding the length of the longest common substring using recursion: // C++ program for finding the length of the longest common substring using recursion #include <iostream> using namespace std; string string_1, string_2; // This is the recursive function that returns the length of the longest common substring int LCS(int iterator_1, int iterator_2, int LCSCount){ // Base Case if (iterator_1 == 0 || iterator_2 == 0) return LCSCount; // If the last character is same then the next character need to be compared to decrement the counter // increase the LCSCount by 1 if (string_1[iterator_1 - 1] == string_2[iterator_2 - 1]) { LCSCount = LCS(iterator_1 - 1, iterator_2 - 1, LCSCount + 1); // if the characters are not the same then you need to decrease the iterator_1 once and update the LCSCount to 0 // The other case can be we can decrement the iterator_2 and update the LCSCount to 0 // The answer is the max of all these cases LCSCount = max(LCSCount, max(LCS(iterator_1, iterator_2 - 1, 0),LCS(iterator_1 - 1, iterator_2, 0))); // return the LCSCount return LCSCount; // Driver code int main(){ int N, M; string_1 = "singer"; string_2 = "singing"; // length of two strings is N = string_1.size(); M = string_2.size(); // call for the recursive function to get the length of the longest common substring cout << LCS(N, M, 0); return 0; Python implementation for finding the length of the longest common substring using recursion: # Python program for finding the length of the longest common substring using recursion # This is the recursive function which returns the length of the longest common substring def LCS(iterator_1, iterator_2, LCSCount): # Base Case if (iterator_1 == 0 or iterator_2 == 0): return LCSCount # If the last character is same then the next character need to be compared to decrement the counter # increase the LCSCount by 1 if (string_1[iterator_1 - 1] == string_2[iterator_2 - 1]): LCSCount = LCS(iterator_1 - 1, iterator_2 - 1, LCSCount + 1) # if the characters are not the same then you need to decrease the iterator_1 once and update the LCSCount to 0 # The other case can be we can decrement the iterator_2 and update the LCSCount to 0 # The answer is the max of all these cases LCSCount = max(LCSCount, max(LCS(iterator_1, iterator_2 - 1, 0), LCS(iterator_1 - 1, iterator_2, 0))) # return the LCSCount return LCSCount # Driver code if __name__ == "__main__": # Given two strings are string_1 = "singer" string_2 = "singing" # length of two strings is N = len(string_1) M = len(string_2) # Call the recursive function by passing the parameters print(LCS(N, M, 0)) Java implementation for finding the length of the longest common substring using recursion: // Java program for finding the length of the longest common substring using recursion class Main{ static String string_1, string_2; // This is the recursive function which returns the length of the longest common substring static int LCS(int iterator_1, int iterator_2, int LCSCount){ // Base Case if (iterator_1 == 0 || iterator_2 == 0){ return LCSCount; // If the last character is same then the next character need to be compared to decrement the counter // increase the LCSCount by 1 if (string_1.charAt(iterator_1 - 1) == string_2.charAt(iterator_2 - 1)){ LCSCount = LCS(iterator_1 - 1, iterator_2 - 1, LCSCount + 1); // if the characters are not the same then you need to decrease the iterator_1 once and update the LCSCount to 0 // The other case can be we can decrement the iterator_2 and update the LCSCount to 0 // The answer is the max of all these cases LCSCount = Math.max(LCSCount, Math.max(LCS(iterator_1, iterator_2 - 1, 0), LCS(iterator_1 - 1, iterator_2, 0))); // return the LCSCount return LCSCount; // Driver code public static void main(String[] args){ // length of two strings is int N, M; // Given two strings are string_1 = "singer"; string_2 = "singing"; N = string_1.length(); M = string_2.length(); // call for the recursive function to get the length of the longest common substring System.out.println(LCS(N, M, 0)); Explanation: The longest common substring from the string_1 and string_2 of lengths N=6 and M=7 respectively, is “sing” and it is of length 4. Therefore the output is 4. Complexity Analysis: Time Complexity: O(3^(N+M)) where N and M are the lengths of the string_1 and string_2 respectively. • Since we are calling the recursive function three times and checking for all characters until the base case is not achieved. • A recursive tree is formed which is: // change the values in the recursive tree put n-1, m-1 generic values. Space Complexity: O(N+M) where N and M are the lengths of the string_1 and string_2 respectively. • Since the recursive function takes a space of N+M for auxiliary stack space during recursion. As we can observe in the recursive tree that the sub-problems are called, again and again, this is the situation of overlapping subproblems which can be resolved by using Dynamic programming. The next approach is explained using Dynamic programming to find the length of the longest common substring. Approach 3: Dynamic programming This approach uses dynamic programming for finding the length of the longest common substring as it is considered to be more efficient in case of overlapping sub-problems for finding the length of the longest common substring. This dynamic programming uses a bottom-up approach to fill up the table. Here, bottom-up approach first focuses on solving the smaller problems at the fundamental level and then integrating them into a whole and giving the complete solution. Dynamic programming is achieved by the tabulation method which uses a 2-D array and stores the length of the longest common substring. This table is filled up in a bottom-up manner. 1. Firstly initialize the 2-D array (DP table) with the size [N+1][M+1]. All the values are initialized to 0. 2. We use one-based indexing as the first row and the first column is 0. 3. This can be done using both iteratively (tabulation) or using memoisation. 4. Iterative approach is discussed below. 5. For filling the values in the table DP. Here dp[i][j] means that the length of the longest substring which ends at ithith character in string_1 and at jthjth character in string_2, and we need to take the maximum of all of these combinations. For filling up the DP table, we must follow the below algorithm: 1. Start two nested loops: 2. First loop with iterator i and second loop with iterator j. 3. If the characters at string_1[i] matches with string_2[j] then the value at dp[i][j] will be equal to the sum of the dp[i-1][j-1] (diagonal value) + 1. ☆ This means that the previous characters of both the strings are also matched and we can increase the LCSCount in the present cell. Therefore the top left diagonal store the count of the length of the longest common substring, until that point. 4. Else the value at dp[i][j] is zero because the string cannot be consecutive i.e, it is not a part of the matching substring. So, we make the present cell value zero. 6. Return the maximum element in the table after formation that is considered to be the length of the longest common substring for the given two strings because each cell represents the length of the longest matching substring till that point. For a clear understanding of this approach, let’s consider an example: Input: string_1 = "acdefb" string_2 = "nbgacde" By following the algorithm mentioned above the DP table obtained is: In the table, we can observe that the row represents string_1 and the column represents string_2. Values in the table are filled accordingly. We can observe that the maximum element in the table is 4, which is the length of the longest common substring which is “acde”. Below is the implementation of this approach: C++ implementation for finding the length of the longest common substring using Dynamic programming: // C++ program for finding the length of the longest substring using Dynamic Programming #include <iostream> #include <string.h> using namespace std; // function which returns the length of longest common substring int LCS(char* string_1, char* string_2, int M, int N){ // This table keeps track of the substrings which are common // and contains the length of the longest substring. // Row in this table represents the string_1 and the column as string_2. // create a 2-D array named dp and intitialise all the values as zero. int dp[M + 1][N + 1]; // To store length of the longest // common substring int LCSCount = 0; // Following steps to fill the values in the dp array in bottom-up manner // Iterative approach for filling the values. // row iterates over the first string for (int i = 0; i <= M; i++){ // coloumn iterates over the second string for (int j = 0; j <= N; j++){ // Base Case if (i == 0 || j == 0) dp[i][j] = 0; // if the characters at string_1 and string_2 matches else if (string_1[i - 1] == string_2[j - 1]){ dp[i][j] = dp[i - 1][j - 1] + 1; // carry forwarding the maximum element in the table LCSCount = max(LCSCount, dp[i][j]); // if the characters did not match then the values is zero dp[i][j] = 0; // return the maximum element in the dp table return LCSCount; // Driver code int main(){ // Given two strings are char string_1[] = "acdefb"; char string_2[] = "nbgacde"; // length of two strings int M = strlen(string_1); int N = strlen(string_2); // call for the LCS function to get the length of the longest common substring cout << LCS(string_1, string_2, M, N); return 0; Python implementation for finding the length of the longest common substring using Dynamic programming: # Python program for finding the length of the longest common substring using dynamic programming def LCS(string_1, string_2, M, N): # This table keeps track of the substrings which are common and contains the length of the longest substring. # Row in this table represents the string_1 and the column as string_2. # create a 2-D array named dp and intitialise all the values as zero. dp = [[0 for k in range(N+1)] for l in range(M+1)] # To store the length of the longest common substring LCSCount = 0 # Following steps to fill the values in the dp array in a bottom-up manner # Iterative approach for filling the values. # row iterates over the first string for i in range(M + 1): # column iterates over the second string for j in range(N + 1): # Base Case if (i == 0 or j == 0): dp[i][j] = 0 # if the characters at string_1 and string_2 matches elif (string_1[i-1] == string_2[j-1]): dp[i][j] = dp[i-1][j-1] + 1 # carry forwarding the maximum element in the table LCSCount = max(LCSCount, dp[i][j]) # if the characters did not match then the values is zero dp[i][j] = 0 # return the maximum element in the dp table return LCSCount # Driver code if __name__ == "__main__": # Given two strings are string_1 = "acdefb" string_2 = "nbgacde" # length of two strings is N = len(string_1) M = len(string_2) # Call the LCS function by passing the parameters print(LCS(string_1, string_2, N, M)) Java implementation for finding the length of the longest common substring using dynamic programming: // java program for finding the length of the longest common substring using dynamic programming class Main { // function which returns the length of longest common substring static int LCS(char string_1[], char string_2[], int M, int N) { // This table keeps track of the substrings which are common // and contains the length of the longest substring. // Row in this table represents the string_1 and the column as string_2. // create a 2-D array named dp and intitialise all the values as zero. int dp[][] = new int[M + 1][N + 1]; // To store length of the longest // common substring int LCSCount = 0; // Following steps to fill the values in the dp array in bottom-up manner // Iterative approach for filling the values. // row iterates over the first string for (int i = 0; i <= M; i++){ // column iterates over the second string for (int j = 0; j <= N; j++){ // # Base Case if (i == 0 || j == 0){ dp[i][j] = 0; // if the characters at string_1 and string_2 matches else if (string_1[i - 1] == string_2[j - 1]){ dp[i][j] = dp[i - 1][j - 1] + 1; // carry forwarding the maximum element in the table LCSCount = Integer.max(LCSCount, dp[i][j]); // if the characters did not match then the values is zero dp[i][j] = 0; // return the maximum element in the dp table return LCSCount; // Driver code public static void main(String[] args){ int M,N; // Given two strings are String string_1 = "acdefb"; String string_2 = "nbgacde"; // length of two strings M = string_1.length(); N = string_2.length(); // call for the LCS function to get the length of the longest common substring System.out.println(LCS(string_1.toCharArray(),string_2.toCharArray(), M, N)); Explanation: The longest common substring from the string_1 and string_2 of lengths N=6 and M=7 respectively, is “sing” and it is of length 4. Therefore the output is 4. Complexity Analysis: Time Complexity: O(M*N) where M and N are the lengths of the given string_1 and string_2 respectively. • For filling up the the DP table (2-D array) in a bottom-up manner, it requires a time of O(N*M). Each cell can be filled as zero or using the previous cell value as required. Space Complexity: O(M*N) where M and N are the lengths of the given string_1 and string_2 respectively. • DP table (2-D array) is of size [M+1][N+1], so the space occupied is O(M*N). Since we are updating the values in the table using the values of the previous updates, space optimization can be done to the above approach. This algorithm is discussed in the below approach. Approach 4: Dynamic Programming- Space Optimization In the above approach, we are only using the last row( updated value) of the 2-D array, therefore we can optimize the space by using a 2-D array of size 2*(min(N, M)). • Algorithm remains the same as filling up of 2-D array but instead of column size as N or M, we instead take 2 because we just need the last row. • Generally, all the dynamic programming approaches can be optimized by using the DP state where we can use the previous state to get the present state’s value and it can also be done using the location of the answer in the DP matrix formed. • However, we need to traverse the two strings using two nested loops, for obtaining the length of the longest substring. • The space optimization is done by : □ if string_1[i – 1] is equal to string_2[j – 1], then dp[i%2][j] = dp[(i – 1)%2][j – 1] + 1 ☆ where dp[i%2][j] represents the present cell and the dp[(i-1)%2][j-1] represents the previous state. ☆ Here we have taken the value to be mod 2, because since it involves only 2 coloumns, we can just use the mod 2 operation while updating the table values. □ Since, we need to carry forward the maximum element in the table, we can do : ☆ LCSCount = max(LCSCount, dp[i%2][j]). □ If the characters are not matched, then the present cell value is zero. • At last, return the length of the longest common substring. Below is the implementation of the space-optimized dynamic programming approach for finding the length of the longest common substring from the given two strings. C++ implementation for finding the length of the longest substring by space optimizing using DP state: // c++ implementation using dynamic programming for space optimisation #include <iostream> #include <string.h> using namespace std; // function which returns the length of longest common substring int LCS(string string_1, string string_2, int M, int N){ // This table keeps track of the substrings which are common // and contains the length of the longest substring. // Row in this table represents the string_1 and the coloum as string_2. // create a 2-D array named dp and intitialise all the values as zero. int dp[2][N + 1]; // To store length of the longest // common substring int LCSCount = 0; // Following steps to fill the values in the dp array in bottom-up manner // Iterative approach for filling the values. for (int i = 1; i <= M; i++){ // column iterates over the second string for (int j = 1; j <= N; j++){ // Base Case if (i == 0 || j == 0) dp[i][j] = 0; // if the characters at string_1 and string_2 matches else if (string_1[i - 1] == string_2[j - 1]){ dp[i%2][j] = dp[(i - 1)%2][j - 1] + 1; // carry forwarding the maximum element in the table LCSCount = max(LCSCount, dp[i%2][j]); // if the characters did not match then the values is zero dp[i%2][j] = 0; // return the maximum element in the dp table return LCSCount; // Driver code int main(){ // Given two strings are string string_1 = "acdefb"; string string_2 = "nbgacde"; // length of two strings int M = string_1.length(); int N = string_2.length(); // call for the LCS function to get the length of the longest common substring cout << LCS(string_1, string_2, M, N); return 0; Python implementation for finding the length of the longest substring by space optimizing using DP state: # Python program for finding the length of the longest common substring using dynamic programming def LCS(string_1, string_2, M, N): # This table keeps track of the substrings which are common and contains the length of the longest substring. # create a 2-D array named dp and intitialise all the values as zero. dp = [[0 for k in range(N+1)] for l in range(2)] # To store the length of longest common substring LCSCount = 0 # Following steps to fill the values in the dp array in bottom-up manner # Iterative approach for filling the values. for i in range(1, M + 1): # column iterates over the second string for j in range(1, N + 1): # Base Case if (i == 0 or j == 0): dp[i][j] = 0 # if the characters at string_1 and string_2 matches elif (string_1[i-1] == string_2[j-1]): dp[i%2][j] = dp[(i-1)%2][j-1] + 1 # carry forwarding the maximum element in the table LCSCount = max(LCSCount, dp[i%2][j]) # if the characters did not match then the values is zero dp[i%2][j] = 0 # return the maximum element in the dp table return LCSCount # Driver code if __name__ == "__main__": # Given two strings are string_1 = "acdefb" string_2 = "nbgacde" # length of two strings is N = len(string_1) M = len(string_2) # Call the LCS function by passing the parameters print(LCS(string_1, string_2, N, M)) Java implementation for finding the length of the longest substring by space optimizing using DP state: // java program for finding the length of the longest common substring using dynamic programming by space optimization technique class Main { // function which returns the length of longest common substring static int LCS(char string_1[], char string_2[], int M, int N) { // This table keeps track of the substrings which are common // and contains the length of the longest substring. // Row in this table represents the string_1 and the coloum as string_2. // create a 2-D array named dp and intitialise all the values as zero. int dp[][] = new int[2][N + 1]; // To store length of the longest // common substring int LCSCount = 0; // Following steps to fill the values in the dp array in bottom-up manner // Iterative approach for filling the values. for (int i = 1; i <= M; i++){ for (int j = 1; j <= N; j++){ // # Base Case if (i == 0 || j == 0){ dp[i][j] = 0; // if the characters at string_1 and string_2 matches else if (string_1[i - 1] == string_2[j - 1]){ dp[i%2][j] = dp[(i - 1)%2][j - 1] + 1; // carry forwarding the maximum element in the table LCSCount = Integer.max(LCSCount, dp[i%2][j]); // if the characters did not match then the values is zero dp[i%2][j] = 0; // return the maximum element in the dp table return LCSCount; // Driver code public static void main(String[] args){ int M,N; // Given two strings are String string_1 = "acdefb"; String string_2 = "nbgacde"; // length of two strings M = string_1.length(); N = string_2.length(); // call for the LCS function to get the length of the longest common substring System.out.println(LCS(string_1.toCharArray(),string_2.toCharArray(), M, N)); Explanation: The longest common substring from the string_1 and string_2 of lengths N=6 and M=7 respectively, is “sing” and it is of length 4. Therefore the output is 4. Complexity Analysis: Time Complexity: O(M*N) where M and N are the lengths of string_1 and string_2 respectively. • For filling up the DP table (2-D array) in a bottom-up manner, it requires a time of O(N*M). Space Complexity: O(min(M*N)) where M and N are the lengths of string_1 and string_2 respectively. • DP table (2-D array) is of size [2][min(N,M)+ 1], so the space occupied is min(M*N)*2. • So the overall space complexity is O(min(M*N)). • We have discussed various approaches for finding the length of the longest substring. • The first approach is by generating all the substrings of the given two strings of length M and N and comparing their substrings. It takes a time complexity of O(N^2 * M^2), this can be optimized by comparing only the same length substrings and it takes the time complexity of O(N^3) where the length of both the strings is the same N. • The second approach uses recursion for finding the length of the longest common substring this takes a time complexity of O(3^(N+M)) and space complexity of O(M+N) • Whenever a recursive approach is present, and a recursive tree is formed we observe that dynamic programming can be used to overcome the re-calculation of the sub-problems again and again. • Instead, we compute the results of all sub-problems and store them. They are extracted whenever required for solving the real problem. • So in this way, we get the third approach which uses dynamic programming for finding the length of the longest common substring from the given two strings. We write DP state and store the values in a 2-D DP table. This takes the time complexity of O(N*M) and space complexity O(N*M) for the DP table. • The fourth approach is the same DP approach but space optimization is done by, simply storing the previous row, therefore the time complexity remains the same i.e, O(N*M) and space complexity is O(2*min(M, N)). HAPPY CODING!!!
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Inversion of parabolic and paraboloidal projections The multidimensional inverse scattering problem for an acoustic medium is considered within the homogeneous background Born approximation. A constant density acoustic medium is probed by a wide-band plane wave source, and the scattered field is observed along a receiver array located outside the medium. The inversion problem is formulated as a generalized tomographic problem. It is shown that the observed scattered field can be appropriately filtered so as to obtain generalized projections of the scattering potential. For a 2-D experimental geometry, these projections are weighted integrals of the scattering potential over regions of parabolic support, whereas they become surface integrals over circular paraboloids for the 2-D case. The inversion problem is therefore similar to that of X-ray tomography, except that instead of being given projections of the object to be reconstructed along straight lines, parabolic or paraboloid projections are given. The inversion procedure that we propose is similar to the X-ray solution, in the sense that it consists of a backprojection operation followed by a 2- or 3-D space invariant filtering. An alternative interpretation of the backprojection operation in terms of a backpropagated field is given. A Projection-Slice Theorem is also derived relating the generalized projections and the scattering potential in the Fourier transform domain. Pub Date: April 1987 □ Acoustic Scattering; □ Born Approximation; □ Homogeneity; □ Inversions; □ Parabolas; □ Arrays; □ Broadband; □ Detection; □ Fourier Transformation; □ Integrals; □ Inverse Scattering; □ Parabolic Bodies; □ Plane Waves; □ Receivers; □ Tomography; □ Weighting Functions; □ X Ray Sources; □ Acoustics
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The discovery that Gauss was most proud of – 07/27/2021 – Marcelo Viana Draw a circle on the paper with a compass. Then, without changing the opening of the compass, draw another circle centered somewhere in the first one. Finally, use a ruler to connect the centers of the two circles to one of the points where they intersect. The figure obtained in this way is an equilateral triangle, that is, the sides of which are all of the same length. The ancient Greeks knew how to use a ruler and compass to build regular polygons with 3, 4, 5 and 15 sides. They also knew how to get from one regular polygon to another with double sides. So they knew how to build the regular hexagon (6 sides) from the equilateral triangle. Can all regular polygons with any number of N sides be constructed with a ruler and compass? The answer is no, but that wasn’t understood until the 18th century when it was proven that regular 7- and 13-sided polygons couldn’t be constructed this way. So what are the constructible values of N, ie so that the regular polygon with N sides can only be constructed with a ruler and compass? The problem caught the attention of none other than the great Carl Friedrich Gauss. In 1796 he showed how to build the regular heptagon (17 pages) with a ruler and compass. This was a discovery Gauss was most proud of. In his great work “Disquisitiones Arithmeticae” he went even further and came to the conclusion that for the construction of a regular polygon it is sufficient that the number N of sides is the product of a power of 2 divided by different Fermat prime numbers. He also stated that this condition would be sufficient, but this was only proven in 1837 by the Frenchman Pierre Pierre de Fermat calculated the numbers of the form 1 plus 2 to 2n for n values from 0 to 4, found them to be prime numbers and held this for all n values. But a few years later Leonhard Euler pointed out that the Fermat number with n = 5 is not a prime number and ironically, until today it has found none other than the five originals he discovered. Since there are 31 different Fermat prime number products, the Gauss-Wantzel theorem gives 31 odd N numbers that can be constructed, and this is the best result known so far. PRESENT LINK: Did you like this column? Subscribers can allow five free hits on any link per day. Just click the blue F below.
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Trigonometry Assignment help homework assignment help is most useful online help portal for the students that providing all Online Trigonometry assignment help Services .Math assignments are common mathematics problems. In most cases, students will have to do math assignments after every math class as homework. Math assignments can be tricky, especially if a student is learning a new mathematics concept. However you can send us Math problems from Basic Math problems like Algebra, Geometry, Trigonometry, and Arithmetic to more advanced Math Assignment Problems like Limits, Continuity, Derivatives, Pre-Calculus, and Double angle formulae The following are important trigonometric relationships (it is unlikely that you will need to know how to prove them and they may be given in your formula book- check!): sin(A + B) = sinAcosB + cosAsinB cos(A + B) = cosAcosB - sinAsinB tan(A + B) = tanA + tanB 1 - tanAtanB To find sin(A - B), cos(A - B) and tan(A - B), just change the + signs in the above identities to - sin(A - B) = sinAcosB - cosAsinB cos(A - B) = cosAcosB + sinAsinB tan(A - B) = tanA - tanB 1 + tanAtanB Double Angle Formulae sin(A + B) = sinAcosB + cosAsinB Replacing B by A in the above formula becomes: sin(2A) = sinAcosA + cosAsinA so sin2A = 2sinAcosA similarly, cos2A = cos²A - sin²A Replacing cos²A by 1 - sin²A (see Pythagorean identities) in the above formula gives: cos2A = 1 - 2sin²A Replacing sin²A by 1 - cos²A gives: cos2A = 2cos²A - 1 It can also be shown that: tan2A = 2tanA 1 - tan²A Pythagorean Identities This important identity can be derived as a direct result of Pythagoras's theorem, when applied to angles in trigonometry: sin²x + cos²x = 1 By dividing each of these terms by sin²x, we can derive a second identity: 1 + cot²x = cosec²x By dividing (1) by cos²x, we arrive at the third (and final) identity: tan²x + 1 = sec²x Radians, like degrees, are a way of measuring angles. One radian is equal to the angle formed when the arc opposite the angle is equal to the radius of the circle. So in the above diagram, the angle ø is equal to one radian since the arc AB is the same length as the radius of the circle. Now, the circumference of the circle is 2pr, where r is the radius of the circle. So the circumference of a circle is 2p larger than its radius. This means that in any circle, there are 2p radians. Therefore 360º = 2p radians. Therefore 180º = p radians. So one radian = 180/p degrees and one degree = p/180 radians. Therefore to convert a certain number of degrees in to radians, multiply the number of degrees by p/180 (for example, 90º = 90 × p/180 radians = p/2). To convert a certain number of radians into degrees, multiply the number of radians by 180/p . Arc Length The length of an arc of a circle is equal to rø, where ø is the angle, in radians, subtended by the arc at the centre of the circle. So in the below diagram, s = rø . Area of Sector The area of a sector of a circle is ½ r² ø, where r is the radius and ø the angle in radians subtended by the arc at the centre of the circle. So in the below diagram, the shaded area is equal to ½ r² ø . For Demo Class Click here Our tutors start working only after the payment is made, to ensure that we are doing work only for serious clients and also our solution meets the required standard. Getting homework help was never so easy you just need to follow following steps: • Send us you Other Assignment or problem through email • Specify the required format such as Word, Excel, Notepad, PDF • Give us a deadline when you need the assignment completed along with the Time Zone. (for example: EST, Australian GMT etc) • Send documents related to your assignment which can help our tutors to provide a better work, any example or format you want the solutions to be in. • Our tutors will review the assignment sent by you and if all the required information is there we will send you the price quoted by our tutor along with the time needed to solve the assignment • You can pay us through paypal or credit card. • After receiving the payment tutors start working on your assignment. • Finally, we deliver the solutions and get a feedback from you regarding our work In case you face any problem or have any query please email us at :- info[@]homeworkassignmenthelp.com
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A triangle has sides A, B, and C. Sides A and B have lengths of 2 and 4, respectively. The angle between A and C is (7pi)/24 and the angle between B and C is (5pi)/8. What is the area of the triangle? | Socratic A triangle has sides A, B, and C. Sides A and B have lengths of 2 and 4, respectively. The angle between A and C is #(7pi)/24# and the angle between B and C is # (5pi)/8#. What is the area of the 1 Answer The area is $\setminus \sqrt{6} - \setminus \sqrt{2}$ square units, about $1.035$. The area is one half the product of two sides times the sine of the angle between them. Here we are given two sides but not the angle between them, we are given the other two angles instead. So first determine the missing angle by noting that the sum of all three angles is $\setminus \ pi$ radians: $\setminus \theta = \setminus \pi - \frac{7 \setminus \pi}{24} - \frac{5 \setminus \pi}{8} = \setminus \frac{\pi}{12}$. Then the area of the triangle is Area $= \left(\frac{1}{2}\right) \left(2\right) \left(4\right) \setminus \sin \left(\setminus \frac{\pi}{12}\right)$. We have to compute $\setminus \sin \left(\setminus \frac{\pi}{12}\right)$. This can be done using the formula for the sine of a difference: $\sin \left(\setminus \frac{\pi}{12}\right) = \setminus \sin \left(\textcolor{b l u e}{\setminus \frac{\pi}{4}} - \textcolor{g o l d}{\setminus \frac{\pi}{6}}\right)$ $= \setminus \sin \left(\textcolor{b l u e}{\setminus \frac{\pi}{4}}\right) \cos \left(\textcolor{g o l d}{\setminus \frac{\pi}{6}}\right) - \setminus \cos \left(\textcolor{b l u e}{\setminus \frac{\ pi}{4}}\right) \sin \left(\textcolor{g o l d}{\setminus \frac{\pi}{6}}\right)$ $= \left(\frac{\setminus \sqrt{2}}{2}\right) \left(\frac{\setminus \sqrt{3}}{2}\right) - \left(\frac{\setminus \sqrt{2}}{2}\right) \left(\frac{1}{2}\right)$ $= \frac{\setminus \sqrt{6} - \setminus \sqrt{2}}{4}$. Then the area is given by: Area $= \left(\frac{1}{2}\right) \left(2\right) \left(4\right) \left(\frac{\setminus \sqrt{6} - \setminus \sqrt{2}}{4}\right)$ $= \setminus \sqrt{6} - \setminus \sqrt{2}$. Impact of this question 1655 views around the world
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Interpolated table lookups using SSE2 [1/2] Something that you are very likely to encounter, when writing SSE2 assembler is the issue of using tables for looking up data. While this is a very simple operation in C, it presents a challenge with a lot of pitfalls, when using SSE2. When to use lookup tables? In general, avoid them, if it is within the possibility. Let’s have a look at a simple example: const float ushort_to_float[65535]; float convert_unsigned_short_to_float(unsigned short value) { return ushort_to_float[value]; This is quite tempting in C – especially if you can do some transformations on the value, by modifying the lookup table, if you need something expensive like the square root, use the inverted value. But lets look at the drawbacks. The table is 256KB, so you will run out of level 1 cache on most CPUs. In C the penalty for a value in L1 cache this is usually a couple of cycles per value, but if it has to be fetched from L2, we are talking about 15 cycles or more on most platforms. But if it saves you a division and a square root or something similar the trade-off may be ok for you. In SSE2, the trade-off picture is changed. Here the penalty per pixel for doing the lookup remains the same (or even rises), while calculating the value is usually only about 25% of cost of what the C-code. So you can see – we can do a lot more calculations in SSE2 before using a table should even be considered. What if I have to use lookup tables? There might be cases, where using tables are the only way, because calculating the values simply aren’t practically possible in SSE2. A practical example of this from Rawstudio is the curve correction, that involves mapping input values to output values based on an n-point 2D curve. Curve Correction: Nearly impossible without using lookup tables. In the “old days”, we used a table with 65535 16 bit entries, that simply gave the output value, much as the example above. The table was 128KB, so again, we had a lot of L1 cache misses. When we re-wrote colour correction, everything internally was converted to float, which made the table 256KB, and even slower. Furthermore it also made sure that we lost a lot of the float precision, by having to fit the float values into 65535 slots. So we got two problems: • Bad precision • Table too large for L1 cache Using lookup tables for float data Since most of these lookups are fairly linear, without too sudden peaks, usually we are fairly in the clear if we reduce the precision of the table, and do linear interpolation between them instead. This avoids posterizing the image data, because all float values will have a unique output. Here is an example, where we use a curve correction with only 256 input/output values, and interpolate them float curve_samples[257]; // Make sure we have an extra entry, so we avoid an // "if (lookup > 255)" later. curve_samples[256] = curve_samples[255]; for (all pixels...) { // Look up and interpolate float lookup = CLAMP(input * 256.0f, 0.0f, 255.9999f); float v0 = curve_samples[(int)lookup]; float v1 = curve_samples[(int)lookup + 1]; lookup -= floorf(lookup); output = v0 * (1.0f - lookup) + v1 * lookup; So now the table is only 1KB in size and fits very nicely within the Level 1 cache. This approach solves both our problems, but at the expense of some calculations. Whether 256 entries are enough is up for you to decide, but the principle applies for all sizes. The “floor” function can be quite slow on some compilers/CPUs, but it is here to demonstrate the algorithm. We will touch this later. In part two, I will look at the actual implementation – there are still a lot of pitfalls in making this work in real life. PS. If you enjoy technical stuff like this, please leave a comment. 7 responses to “Interpolated table lookups using SSE2 [1/2]” 1. I can’t say I know too much about the technical details of image processing, but I do appreciate the insight post like these give. Thanks! It seems like you could reduce the size further is you used a different data type for the LUT, for instance a 16-bit or 8-bit int. With my monitor at 1024 pixels high, I couldn’t imagine being able to enter a curve at more than 10-bits precision, and that’s if the screen is maximized. Sounds like the image data is in float, but does it take too many cycles to cast an int to a float? I wonder if the loss of precision would be noticable. Anyways, these are just questions floating in my head, and I’m not trying to criticize your implementation, you folks know way more than I ever will about this stuff. □ The game is a little different when you are dealing with Raw data, because you operate in a linear colorspace. Because of this you need additional precision, since errors become very visible in dark areas. See this PDF on why you need at least 16 bit precision for Raw images: 2. Yes, I enjoyed this posting very much. Thanks! 3. Thanks, I enjoy this kind of blog posts! Keep them coming :) 4. […] Continued from part 1… […] 5. The problem with that approach is that decreasing the table precision leads to at least two problems. First, the curve downscaler 65536->256 has to be chosen carefully. And then such a big downscale can’t be reasonably inversed using a simple linear interpolation w/o introducing aliasing problems. Do you have some PSNR numbers comparing the original results using interpolation between the 65536 sampling points of the curve, and the faster 256 entry table interpolation ? □ Hi Edouard! >First, the curve downscaler 65536->256 has to be chosen carefully. The curve isn’t downscaled as such – the curve is simply sampled with 256 entries instead of 65536. So we get 256 exact values, and not some that are derived from 65536. >And then such a big downscale can’t be reasonably inversed using a simple linear interpolation w/o introducing aliasing problems. How do you get to that result? The curve is an input->output mapping, where neutral is input=output (straight line). You would have to have a curve that has variations where the correction is changed at a 1/128th pixel resolution, which simply wouldn’t make sense. So unless you have information with a frequency close to or above the nyquist frequency, a linear interpolation should not generate aliasing. If you consider the image of the curve, you have approximately the horizontal resolution of the image (the image is ~290 pixels wide). The vertical sample point have float point precision, so the vertical resolution is very precise. The vertical resolution is the reason there is no precision issues, even with such small tables. >Do you have some PSNR numbers comparing the original results using interpolation between the 65536 sampling points of the curve, and the faster 256 entry table interpolation ? I’ll make a few calculations that calculates the min, max and average error on “realistic” and some “extreme” curves. Klaus Post on October 28, 2010
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. Calculate the median. - Datapott Analytics For the Week 1 Complete assignment, solve each problem and answer each question that corresponds with it. Explain how you arrived at the answer for each problem. The total word count for your assignment should be a minimum of 1200 words. Identify the type of data (qualitative/quantitative) and the level of measurement Explain your choice. The average monthly temperature in degrees Fahrenheit for the city of Wilmington, Delaware, through-out the year. The ages of the respondents in a survey. The years in which the respondents to a survey were born. The ethnicity of the respondents in a survey classi-fied as White, African American, Asian, or Other. The Arizona Diamondbacks and the city of Phoenix are considering building a new baseball stadium for the Major League Baseball team. A decision needs to be made whether to enclose the new facility with a roof. To help make this decision, data were gathered on the number of days per month it rained during the baseball season, which are shown in the following table. a. Construct a frequency distribution for these data. b. Using the results from part a, calculate the rela-tive frequencies for each class. c. Using the results from part a, calculate the cumulative relative frequencies for each class. d. Construct a histogram for these data. e. What is the likelihood that a month during the baseball season will experience one day or less The following table lists the math SAT scores for 72 students. a. Using the 2k ≥ n rule, construct a frequency distribution for these data. b. Using the results from part a, calculate the relative frequencies for each class. c. Using the results from part a, calculate the cumulative relative frequencies for each class. d. Construct a histogram for these data. 7. The following data represent the number of days homes were on the market before being sold in New Castle County, Delaware a. Calculate the mean. b. Calculate the median. c. Determine the mode. d. Which of these three measurements would best describe the central tendency of the data? e. Describe the shape of this 8. Internet service providers compete on services such as download speeds that are measured in megabits per second (Mbps). The following data represent the download speeds that 16 Comcast customers experienced recently. a. Calculate the mean. b. Calculate the median. c. Determine the mode. d. Describe the shape of the distribution
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how to make types of angles math’s working model(TLM project) - Science Projects | Maths TLM | English TLM | Physics Projects | Computer Projects | Geography Projects | Chemistry Projects | Working Projects | Working Models | DIY for School / College Science Exhibitions or Fair how to make types of angles math’s working model(TLM project) In this post we written on how to make the model on types of angles math’s working project (TLM project) using cardboard and color paper types of angles math’s working model(TLM project) Creating a working model to demonstrate types of angles using cardboard and color paper is an excellent way to visualize and understand different angle measurements. Here’s a step-by-step guide to making the model: Materials you will need: • Cardboard (for the base and angle representations) • Color paper (different colors for different types of angles) • Scissors • Glue or adhesive • Markers or pens (for adding details) Step-by-step instructions: 1. Prepare the base: □ Take a large piece of cardboard to serve as the base for your angle model. 2. Draw and cut out angle representations: □ Use a ruler and a pencil to draw various angle representations on the color paper. For example: ☆ Right angle (90 degrees): Draw an “L” shape with one angle measuring 90 degrees. ☆ Acute angle (less than 90 degrees): Draw a small angle less than 90 degrees. ☆ Obtuse angle (greater than 90 degrees): Draw a larger angle greater than 90 degrees. ☆ Straight angle (180 degrees): Draw a straight line. □ Cut out these angle representations from the color paper using scissors. 3. Glue the angles onto the cardboard: □ Glue each angle representation onto the cardboard base, leaving some space between them. 4. Demonstrate angle relationships: □ Use the model to demonstrate the relationships between different angle types. For example, show how two acute angles can add up to form a right angle, or how two supplementary angles form a straight angle. This model provides a visual representation of various angle measurements and their characteristics. It’s a fun and educational project to understand the different types of angles and how they relate to one another in geometry. #typesofangles #mathsworkingmodel #tlm #craftpiller #maths #workingmodel #diy Video guide on making types of angles math’s working model(TLM project) Leave a Comment
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The Remainders Game If you haven't already seen Remainders, it would be worth trying that task before playing this game. The computer will think of a number between 1 and 100. Can you work out what it is? Choose a divisor and the computer will give you some information about the number. The fewer divisions you require, the more points you get. If you identify the number correctly after 1 or 2 divisions - you gain 12 points 3 divisions - you gain 11 points 4 divisions - you gain 9 points 5 divisions - you gain 6 points 6 divisions - you gain 3 points 7 or more divisions - you gain 1 point If you guess wrongly you lose 15 points, even if your guess satisfies all the criteria. How soon can you reach 100 points? Challenge your friends to do better. The modulator will help you play the game, but eventually, try to play the game without it. Click on the purple cog to open the settings menu and change the level of the game: In Level 1 the chosen number will be from 1 to 60 inclusive and you can divide by each of the numbers from 1 to 10. In Level 2 (the default setting), the chosen number will be from 1 to 100 inclusive and you can divide by each of the numbers from 2 to 10. In Level 3 the chosen number will be from 1 to 100 inclusive and you can only divide by a selection of the numbers from 1 to 10. This page contains some examples to clarify the rules for the Remainders Game. If you guess wrong you lose 15 points, so don't guess until you are certain that there is only one possible answer. Suppose you have discovered that your number has a remainder of $1$ when divided by $2$, and a remainder of $0$ when divided by $5$ At this point, you know that the hidden number is both odd and a multiple of $5$. You might think of $15$ as a number that would fit with both of these criteria. However, there are still other possibilities that the number could be, such as $5$, $25$ or even $95$. This means that you'd need to do some more divisions to work out which number it is. Try to ensure that each division you carry out provides new information - it should rule out some numbers. Suppose you have divided by $4$ and found that you are left with a remainder of $1$. At this stage, you know your number is $1$ or $5$ or $9$ or..., so all the possibilities that you get are odd. Therefore, we already know that, if you choose to divide by $2$, the remainder will be So, if you did divide by $2$ in this situation you would gain no extra information as to what the hidden number is. If you carry out unnecessary divisions your score for getting the right answer will be lower. You can use The Modulator tool to help you to decide whether a division will rule out any numbers.
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ThmDex – An index of mathematical definitions, results, and conjectures. This result is a subresult to R5185: Tight lower bound to a finite product of positive real numbers . In that result, we see that $\prod_{n = 1}^N x_n$ attains the lowest possible value when $x_1 = x_2 = \cdots = x_N$. In this case, we have \begin{equation} \left( \frac{1}{\frac{1}{N} \sum_{n = 1}^ N \frac{1}{x_n}} \right)^N = \prod_{n = 1}^N x_n = x^N_k \end{equation} for some $k = 1, \ldots, N$. Substituting $\sum_{n = 1}^N \frac{1}{x_n} = a$ and raising both sides to power $1 / N$, we have \ begin{equation} \frac{N}{a} = \frac{1}{a / N} = x_k \end{equation} Since $k = 1, \ldots, N$ was arbitrary, the claim follows. $\square$
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Live, Online, Honors Level, Homeschool Math Classes for Pre-Calculus This class is designed to help homeschooling families succeed in their study of pre-calculus topics including functions, analytical geometry, probability, series, sequences, and trigonometry by providing quality instruction and objective evaluation using Foerster’s Algebra and Trigonometry: Functions and Applications (Pearson Prentice Hall: 2005). This rigorous and highly regarded textbook is designed to be completed in two years, so this class covers chapters 9 – 15. Students will be expected to spend roughly an hour and a half to two hours on homework daily. Class will meet two days per week for one hour and five minutes each session for instruction. Since this is a homeschooling tutorial, parents will have the certain responsibilities as stated here. Understand that this is an honors level class and not designed for the struggling or lazy math student. Students will be expected to complete homework assignments before the subsequent class, participate in class with a working microphone, and proactively seek help when not understanding the material. Students need to have completed Geometry as well as Algebra II using Foerster’s text with an A or B average or gain special permission from the instructor. It is recommended that students be at least 15 years old, entering at least 11th grade. Parents of younger students should contact the instructor before registering for the class. The text used has recently changed publishers, but under the cover, the editions are exactly the same. *Students may want to use a Graph Ruled Spiral Notebook for all of their work so they can draw graphs beside their calculations. “Thank you, Mrs. Flynn, for the excellent preparation for AP Calculus that our daughter received in your Advanced Math (now named Pre-Calculus) class. After your online course, not only did she successfully complete the AP class, but she took the AP Calculus AB/BC exam and achieved a 5. This gave her a great edge on her college applications and she has been accepted into the Honors Fellows program at her top choice college. Thank you for your part in her math success!” We would love to sign up [our daughter] for your Advanced Math class. I so appreciate what you have done wit…Read more
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September 2015 - SignalSolver This is a Berkshire Hathaway trading strategy which would have given almost ten times the return performance of buy/hold over the last 10 years with half the drawdown. The strategy is detailed in the table below, it was straightforward, with 123 trades over the 10 year backtest period. All trading was done at the weekly close of business. This was a strategy where there was usually a buy and a sell signal every week (398 out of 528 weeks), but selling was initiated by the presence of a sell signal and an absence of a buy signal. There was a strong buy bias, appropriate for the underlying positive trend of the stock. Equity curve, signals and positions are shown in the growth chart below. Notice the preponderance of white signals which indicate dual signal days. If the sell point was set at zero percent, the algorithm gave positive results for all buy points greater than 0.25%. Below we show how return varies with buy and sell parameters. The returns for every combination of buy and sell parameter are shown in the surface plot below. Looking at the minimum annualized returns for each of the four 132 week periods, you can see that the algorithm was much better behaved than the underlying stock, worst case was 19.26% which happened in the most recent quartus. The long side of the algorithm showed a min return of 15% with a stddev of only 3.37%, which is quite tight. As always, this is not a recommendation to trade using this algorithm, just an interesting backtest result. For a list of trades, see here: BRKA.W Trades. Please note, the above result was corrected 12/28/2015 to address a bug fix in the short side calculations. Update Oct 21st 2016 GOOGL Trading Strategy (Weekly) Here is a Google Inc (GOOGL) trading strategy with once a week intervention which would have performed significantly better than buy-hold over the last 10 years. Annualized return was 31.5% vs. 16.5% (returning $149K for 10K outlay vs $36.7K, compounded), drawdown was 40.2% vs. 62.4%, so reward/risk was better. The buy and cover signal (see table below) was present every week where the price dropped 0.03% below the last price the stock was sold at which happened 174 times over the course of the 528 weeks of the analysis, and usually happened the week following a sell/short. The sell and short signal happened every week the stock price rose 8.13% or more above the open price of the current week. This happened only 29 times, so there is a strong buy-side bias to this strategy. All trading would have been done at the open of the week following the signal. The equity curve for this GOOGL trading strategy shows that the stock was held short (the red bands in the background) for small periods, typically a week. The scan below shows how the annualized return changed, had the buy and sell parameters changed. At 2.35% buy point, the algorithm gets stuck, resulting in a loss. This is characteristic of trading strategies which reference buy or sell prices. At 1.5% sell point and below, the algorithm made a loss, but returns for all buy points above that were positive, for the 0.03% buy point. One nice characteristic this algorithm had was consistent returns for each of the 132 week periods in the backtest. You can see from the table below that the annualized return was between 28% and 34.8% for every period. Compare that with buy-hold which ranged from -4.6% to 26% You can also see this characteristic on the scans for each quartus: You can download the list of trades in .xlsx format: GOOGL.W Trades. As of Sept 16th 2015, the algorithm is long, awaiting a sell signal if the price hits 708.933. Last sell price was 654.34. Please note, while this is an interesting backtest result, it is not a suggestion to trade this way. As always I don’t know how this strategy will fare in the future, but will track it from time to This post was edited 12/28/15 to correct an error in the short-side returns. Update Oct 21st 2016 This strategy has pretty much followed buy-hold long: DIS Trading Strategy Here is a Walt Disney Company (DIS) trading strategy with daily maintenance which had quite nice characteristics; 59% annualized return over the last 2 years. The performance was better than long buy-hold with lower drawdown and about three times the reward-risk. Signal reinforcement was good, and not many dual signal days. Below I show the return for the DIS trading strategy for each 6 month period going back 2 years. Comparing with long buy-hold you can see better consistency (higher minimum and lower StdDev). Please note that this is simply a measurement, not an opinion or financial advice. This post was corrected 1/8/2016 to correct a miscalculation in the short-side returns. Update Oct 21st 2016 Strategy peaked 11/23/15, shortly after the stock price peaked. Two GLD Trading Strategies (Daily) GLD is the much traded SPDR Gold Trust ETF. I find these two GLD trading strategies interesting because they gave reasonable results (32.6% and 48% annualized return) for each of the four 6 month periods of the analysis. The strategies require daily intervention. Strategy 1: BCS AHC This is a buy on fall, sell on rise strategy using the close price as the buy reference and the high price as the sell reference. As you can see from the life chart and the longevity analysis, the lowest return for the last four 6 month periods was close to 30% annualized. You can easily find algorithms with over 50% annualized return for GLD, but they are not as consistent, with lowest quartus returns of around 14%. I would prefer to see more reinforcement on the signals, but there it is. As of Sat Sept 5th, this strategy is Short with no transactions pending. You can view the trades in spreadsheet format here: GLD.D Trades Strategy 2: AOO AHCI In many ways this strategy shows better results than the BCS AHC strategy, for example there was lower drawdown, higher return, better signal reinforcement and good consistency (minimum 6 month quartus return of 38.43%). On the other hand, signals were cluttered, with 67 dual signal days, 50 buy signal days and 75 sell signal days. Also, its not an trading strategy that makes intuitive sense; buy on rise, sell on rise. Maybe it is one of those serendipitous occurrences, we shall see. Notice the sell strategy is almost the same as for BCS AHC but the sell signal percentage is quite As of Sat Sept 5th, the strategy is Short with no transactions pending. For a list of trades in Excel format: GLD.D2 Trades. For a more detailed explanation of the above charts, please go here. Algorithms were discovered by SignalSolver. Please note, the above analysis was corrected on 12/28/2015 to reflect a bug fix in SignalSolver. Original returns were $9530 and $12753 respectively. Update Oct 21st 2016 Both algorithms peaked 12/30/2015.
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libquadmath/printf/mul.c - gcc - Git at Google /* mpn_mul -- Multiply two natural numbers. Copyright (C) 1991, 1993, 1994, 1996 Free Software Foundation, Inc. This file is part of the GNU MP Library. The GNU MP Library is free software; you can redistribute it and/or modify it under the terms of the GNU Lesser General Public License as published by the Free Software Foundation; either version 2.1 of the License, or (at your option) any later version. The GNU MP Library is distributed in the hope that it will be useful, but WITHOUT ANY WARRANTY; without even the implied warranty of MERCHANTABILITY or FITNESS FOR A PARTICULAR PURPOSE. See the GNU Lesser General Public License for more details. You should have received a copy of the GNU Lesser General Public License along with the GNU MP Library; see the file COPYING.LIB. If not, write to the Free Software Foundation, Inc., 59 Temple Place - Suite 330, Boston, MA 02111-1307, USA. */ #include <config.h> #include "gmp-impl.h" /* Multiply the natural numbers u (pointed to by UP, with USIZE limbs) and v (pointed to by VP, with VSIZE limbs), and store the result at PRODP. USIZE + VSIZE limbs are always stored, but if the input operands are normalized. Return the most significant limb of the NOTE: The space pointed to by PRODP is overwritten before finished with U and V, so overlap is an error. Argument constraints: 1. USIZE >= VSIZE. 2. PRODP != UP and PRODP != VP, i.e. the destination must be distinct from the multiplier and the multiplicand. */ /* If KARATSUBA_THRESHOLD is not already defined, define it to a value which is good on most machines. */ #ifndef KARATSUBA_THRESHOLD #define KARATSUBA_THRESHOLD 32 #if __STDC__ mpn_mul (mp_ptr prodp, mp_srcptr up, mp_size_t usize, mp_srcptr vp, mp_size_t vsize) mpn_mul (prodp, up, usize, vp, vsize) mp_ptr prodp; mp_srcptr up; mp_size_t usize; mp_srcptr vp; mp_size_t vsize; mp_ptr prod_endp = prodp + usize + vsize - 1; mp_limb_t cy; mp_ptr tspace; if (vsize < KARATSUBA_THRESHOLD) /* Handle simple cases with traditional multiplication. This is the most critical code of the entire function. All multiplies rely on this, both small and huge. Small ones arrive here immediately. Huge ones arrive here as this is the base case for Karatsuba's recursive algorithm below. */ mp_size_t i; mp_limb_t cy_limb; mp_limb_t v_limb; if (vsize == 0) return 0; /* Multiply by the first limb in V separately, as the result can be stored (not added) to PROD. We also avoid a loop for zeroing. */ v_limb = vp[0]; if (v_limb <= 1) if (v_limb == 1) MPN_COPY (prodp, up, usize); MPN_ZERO (prodp, usize); cy_limb = 0; cy_limb = mpn_mul_1 (prodp, up, usize, v_limb); prodp[usize] = cy_limb; /* For each iteration in the outer loop, multiply one limb from U with one limb from V, and add it to PROD. */ for (i = 1; i < vsize; i++) v_limb = vp[i]; if (v_limb <= 1) cy_limb = 0; if (v_limb == 1) cy_limb = mpn_add_n (prodp, prodp, up, usize); cy_limb = mpn_addmul_1 (prodp, up, usize, v_limb); prodp[usize] = cy_limb; return cy_limb; tspace = (mp_ptr) alloca (2 * vsize * BYTES_PER_MP_LIMB); MPN_MUL_N_RECURSE (prodp, up, vp, vsize, tspace); prodp += vsize; up += vsize; usize -= vsize; if (usize >= vsize) mp_ptr tp = (mp_ptr) alloca (2 * vsize * BYTES_PER_MP_LIMB); MPN_MUL_N_RECURSE (tp, up, vp, vsize, tspace); cy = mpn_add_n (prodp, prodp, tp, vsize); mpn_add_1 (prodp + vsize, tp + vsize, vsize, cy); prodp += vsize; up += vsize; usize -= vsize; while (usize >= vsize); /* True: usize < vsize. */ /* Make life simple: Recurse. */ if (usize != 0) mpn_mul (tspace, vp, vsize, up, usize); cy = mpn_add_n (prodp, prodp, tspace, vsize); mpn_add_1 (prodp + vsize, tspace + vsize, usize, cy); return *prod_endp;
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Math Contest Repository Hypatia 2024 Question 2, CEMC UWaterloo (Hypatia 2024, Question 2, CEMC - UWaterloo) For a positive $3$-digit integer $n$, $f(n)$ is equal to the sum of $n$ and the digits of $n$. For example, $f(351)=351+3+5+1=360$. Note: The decimal representation of a $3$-digit number $abc$ is $a \cdot 10^2 + b \cdot 10 + c$. For example, $836 = 8 \cdot 10^2 + 3 \cdot 10 + 6$. (a) What is the value of $f(132)$? (b) If $f(n)=175$, what is the value of $n$? (c) If $f(n)=204$, determine all possible values of $n$. Answer Submission Note(s) In part (c), sort your answers in ascending order and separate them with a comma. Separate the answers for each part with a single space. For example: "a b c1,c2" Please login or sign up to submit and check if your answer is correct. flag Report Content You should report content if: • It may be offensive. • There is something wrong with it (statement or difficulty value) • It isn't original. Thanks for keeping the Math Contest Repository a clean and safe environment!
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itertools.combinations_with_replacement() python | HackerRank Solutions itertools.combinations_with_replacement() python Problem Statement : itertools.combinations_with_replacement(iterable, r) This tool returns r length subsequences of elements from the input iterable allowing individual elements to be repeated more than once. Combinations are emitted in lexicographic sorted order. So, if the input iterable is sorted, the combination tuples will be produced in sorted order. Sample Code >>> from itertools import combinations_with_replacement >>> print list(combinations_with_replacement('12345',2)) [('1', '1'), ('1', '2'), ('1', '3'), ('1', '4'), ('1', '5'), ('2', '2'), ('2', '3'), ('2', '4'), ('2', '5'), ('3', '3'), ('3', '4'), ('3', '5'), ('4', '4'), ('4', '5'), ('5', '5')] >>> A = [1,1,3,3,3] >>> print list(combinations(A,2)) [(1, 1), (1, 3), (1, 3), (1, 3), (1, 3), (1, 3), (1, 3), (3, 3), (3, 3), (3, 3)] You are given a string S. Your task is to print all possible size k replacement combinations of the string in lexicographic sorted order. Input Format A single line containing the string S and integer value k separated by a space. The string contains only UPPERCASE characters. Output Format Print the combinations with their replacements of string S on separate lines. Solution : Solution in C : from itertools import combinations_with_replacement s,k = input().split(' ') k = int(k) s = [x for x in s] l = list(combinations_with_replacement(s,k)) for i in l: s = '' for k in i: View More Similar Problems Left Rotation A left rotation operation on an array of size n shifts each of the array's elements 1 unit to the left. Given an integer, d, rotate the array that many steps left and return the result. Example: d=2 arr=[1,2,3,4,5] After 2 rotations, arr'=[3,4,5,1,2]. Function Description: Complete the rotateLeft function in the editor below. rotateLeft has the following parameters: 1. int d View Solution → Sparse Arrays There is a collection of input strings and a collection of query strings. For each query string, determine how many times it occurs in the list of input strings. Return an array of the results. Example: strings=['ab', 'ab', 'abc'] queries=['ab', 'abc', 'bc'] There are instances of 'ab', 1 of 'abc' and 0 of 'bc'. For each query, add an element to the return array, results=[2,1,0]. Fun View Solution → Array Manipulation Starting with a 1-indexed array of zeros and a list of operations, for each operation add a value to each of the array element between two given indices, inclusive. Once all operations have been performed, return the maximum value in the array. Example: n=10 queries=[[1,5,3], [4,8,7], [6,9,1]] Queries are interpreted as follows: a b k 1 5 3 4 8 7 6 9 1 Add the valu View Solution → Print the Elements of a Linked List This is an to practice traversing a linked list. Given a pointer to the head node of a linked list, print each node's data element, one per line. If the head pointer is null (indicating the list is empty), there is nothing to print. Function Description: Complete the printLinkedList function in the editor below. printLinkedList has the following parameter(s): 1.SinglyLinkedListNode View Solution → Insert a Node at the Tail of a Linked List You are given the pointer to the head node of a linked list and an integer to add to the list. Create a new node with the given integer. Insert this node at the tail of the linked list and return the head node of the linked list formed after inserting this new node. The given head pointer may be null, meaning that the initial list is empty. Input Format: You have to complete the SinglyLink View Solution → Insert a Node at the head of a Linked List Given a pointer to the head of a linked list, insert a new node before the head. The next value in the new node should point to head and the data value should be replaced with a given value. Return a reference to the new head of the list. The head pointer given may be null meaning that the initial list is empty. Function Description: Complete the function insertNodeAtHead in the editor below View Solution →
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Source term bricks (and Neumann condition) Source term bricks (and Neumann condition)¶ This brick adds a source term, i.e. a term which occurs only in the right hand side of the linear (tangent) system build by the model. If \(f\) denotes the value of the source term, the weak form of such a term is \[\int_{\Omega} f v\ dx\] where \(v\) is the test function. The value \(f\) can be constant or described on a finite element method. It can also represent a Neumann condition if it is applied on a boundary of the domain. The function to add a source term to a model is: add_source_term_brick(md, mim, varname, dataexpr, region = -1, directdataname = std::string()); where md``is the model object, ``mim is the integration method, varname is the variable of the model for which the source term is added, dataexpr has to be a regular expression of GWFL, the generic weak form language (except for the complex version where it has to be a declared data of the model). It has to be scalar or vector valued depending on the fact that the variable is scalar or vector valued itself. region is a mesh region on which the term is added. If the region corresponds to a boundary, the source term will represent a Neumann condition. directdataname is an optional additional data which will directly be added to the right hand side without assembly. The brick has a working complex version. A slightly different brick, especially dedicated to deal with a Neumann condition, is added by the following function: add_normal_source_term_brick(md, mim, varname, dataexpr, region); The difference compared to the basic source term brick is that the data should be a vector field (a matrix field if the variable varname is itself vector valued) and a scalar product with the outward unit normal is performed on it.
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Tree transducer explained In theoretical computer science and formal language theory, a tree transducer (TT) is an abstract machine taking as input a tree, and generating output – generally other trees, but models producing words or other structures exist. Roughly speaking, tree transducers extend tree automata in the same way that word transducers extend word automata. Manipulating tree structures instead of words enable TT to model syntax-directed transformations of formal or natural languages. However, TT are not as well-behaved as their word counterparts in terms of algorithmic complexity, closure properties, etcetera. In particular, most of the main classes are not closed under composition. The main classes of tree transducers are: Top-Down Tree Transducers (TOP) A TOP T is a tuple such that: • is a finite set, the set of states; • is a finite ranked alphabet, called the input alphabet; • is a finite ranked alphabet, called the output alphabet; • is a subset of Q, the set of initial states; and • is a set of rules of the form , where is a symbol of Σ, is the arity of is a state, and is a tree on Γ and , such pairs being nullary. Examples of rules and intuitions on semantics For instance, is a rule – one customarily writes instead of the pair – and its intuitive semantics is that, under the action of , a tree with at the root and three children is transformed into where, recursively, are replaced, respectively, with the application of on the first child and with the application of on the third. The semantics of each state of the transducer T, and of T itself, is a binary relation between input trees (on Σ) and output trees (on Γ). A way of defining the semantics formally is to see as a term rewriting system, provided that in the right-hand sides the calls are written in the form , where states are unary symbols. Then the semantics of a state is given by The semantics of is then defined as the union of the semantics of its initial states: Determinism and domain As with tree automata, a TOP is said to be deterministic (abbreviated DTOP) if no two rules of δ share the same left-hand side, and there is at most one initial state. In that case, the semantics of the DTOP is a partial function from input trees (on Σ) to output trees (on Γ), as are the semantics of each of the DTOP's states. The domain of a transducer is the domain of its semantics. Likewise, the image of a transducer is the image of its semantics. Properties of DTOP • DTOP are not closed under union: this is already the case for deterministic word transducers. • The domain of a DTOP is a regular tree language. Furthermore, the domain is recognisable by a deterministic top-down tree automaton (DTTA) of size at most exponential in that of the initial DTOP. That the domain is DTTA-recognizable is not surprising, considering that the left-hand sides of DTOP rules are the same as for DTTA. As for the reason for the exponential explosion in the worst case (that does not exist in the word case), consider the rule . In order for the computation to succeed, it must succeed for both children. That means that the right child must be in the domain of . As for the left child, it must be in the domain of . Generally, since subtrees can be copied, a single subtree can be evaluated by multiple states during a run, despite the determinism, and unlike DTTA. Thus the construction of the DTTA recognising the domain of a DTOP must account for of states and compute the intersections of their domains, hence the exponential. In the special case of DTOP, that is to say DTOP where each appears at most once in the right-hand side of each rule, the construction is linear in time and space. • The image of a DTOP is not a regular tree language. Consider the transducer coding the transformation ; that is, duplicate the child of the input. This is easily done by a rule , where encodes the . Then, absent any restrictions on the first child of the input, the image is a classical non-regular tree language. • However, the domain of a DTOP cannot be restricted to a regular tree language. That is to say, given a DTOP T and a language L, one cannot in general build a DTOP such that the semantics of is that of This property is linked to the reason deterministic top-down tree automata are less expressive than bottom-up automata: once you go down a given path, information from other paths is inaccessible. Consider the transducer coding the transformation ; that is, output the right child of the input. This is easily done by a rule , where encodes the identity. Now let's say we want to restrict this transducer to the finite (and thus, in particular, regular) domain . We must use the rules . But in the first rule, does not appear at all, since nothing is produced from the left child. Thus, it is not possible to test that the left child is . In contrast, since we produce from the right child, we can test that it is . In general, the criterion is that DTOP cannot test properties of subtrees from which they do not produce output. • DTOP are not closed under composition. However this problem can be solved by the addition of a lookahead: a tree automaton, coupled to the transducer, that can perform tests on the domain which the transducer is incapable of.^[2] This follows from the point about domain restriction: composing the DTOP encoding identity on with the one encoding must yield a transducer with the semantics , which we know is not expressible by a DTOP. • The typechecking problem—testing whether the image of a regular tree language is included in another regular tree language—is decidable. • The equivalence problem—testing whether two DTOP define the same functions—is decidable.^[3] Bottom-Up Tree Transducers (BOT) As in the simpler case of tree automata, bottom-up tree transducers are defined similarly to their top-down counterparts, but proceed from the leaves of the tree to the root, instead of from the root to the leaves. Thus the main difference is in the form of the rules, which are of the form • Book: Hubert. Comon. Max. Dauchet. Rémi. Gilleron. Florent. Jacquemard. Denis. Lugiez. Christof. Löding. Sophie. Tison. Marc. Tommasi. Tree Automata Techniques and Applications. November 2008. Chapter 6: Tree Transducers . https://hal.inria.fr/hal-03367725/document. 11 February 2014. . • Book: Hosoya, Haruo. Cambridge University Press. 978-1-139-49236-2. Foundations of XML Processing: The Tree-Automata Approach. 4 November 2010. Notes and References 1. Baker, B.S.: Composition of top-down and bottom-up tree transductions. Inf. Control 41(2), 186–213 (1979) 2. Maneth. Sebastian. A Survey on Decidable Equivalence Problems for Tree Transducers. International Journal of Foundations of Computer Science. December 2015. 26. 8. 1069–1100. 10.1142/ S0129054115400134. 20.500.11820/2f1acef4-1b06-485f-bfd1-88636c9e2fe6. free. 3. Decidability results concerning tree transducers I. www.inf.u-szeged.hu.
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Objections to the Kalam Cosmological Argument for God The Kalaam faces many issues. There are two big issues for me. First, because of premise 2, the argument is tied heavily with science and scientific discovery. This makes it dependent of scientific trends and to the winds of change. The cutting edge of scientific discussion will change the various aspects of causality, but also the nature of time itself. In other words, unless you constantly update the thing and hope the science continues in your favor, you are going to face a problem. Further, this makes it pragmatically useless to discuss in most cases because you need to understand the science behind it. Without the science, you effectively are just appealing to whatever authority you like. Second, and probably most importantly, the argument must turn into a different argument in a very short order. You will need to go back and defend another cosmological argument or similar proof of god because the premises are highly questionable. The argument takes an observation about what occurs in our universe- "Things that begin to exist have a cause,"- and then applies this observation to the nature of reality outside of the universe. It is currently unknown how reality outside of our universe would operate, if it exists at all. I could go on, but I think there are enough criticisms of the argument that stick to make it problematic. Chany OptionsShare Need some counter-examples of causal closure. every physical effect (i.e. caused event) has physical sufficient causes — http://philpapers.org/archive/VICOTC.pdf Fishfry and, I expect most of the others that have responded to your thread, is/are perfectly familiar with Hilbert's Hotel, and why it is not an argument for any statement other than 'aren't infinities interesting?' Ditto for Aristotle's notions of potential and actual infinities. — andrewk (Y) Hilbert’s Hotel Shandy’s Diary , for example, are what we call veridical paradoxes, and do not imply a contradiction, but they do show some counter-intuitive implications of infinites. Some select objections from the trenches … • even if sound, the argument does not suggest anything “divine”, sentient, conscious, thinking, caring, loving, warranting worship or prayer, so the argument requires more to be of particular • if gods/God can be atemporal (changeless, “outside of time”, or something), assuming that makes sense, then we might suppose any such “origin” of the universe • if there was a definite earliest time (or “time zero”), then anything that existed at that time, began to exist at that time, and that includes any first causes, gods/God, or whatever else • phrases like “before time” and “a cause of causation” are incoherent, violates identity (the 1^st law) • by contemporary cosmology (e.g. Big Bang) spacetime is an aspect of the universe … □ spatiality and temporality are aspects of the universe □ causation is temporal □ therefore causation is an aspect of the universe, not somehow “outside” • dichotomistically … □ if some God of theism could create something out of “nothing”, as it were, then nihil fit ex nihilo is already violated, and we might as well dispose of the principle, in which case said God is an extraneous hypothesis □ if some God of theism created the universe from something already existing, then whatever comprise the universe “always” existed, perhaps “eternally” (to the extent that’s meaningful), and we might as well dispose of the extras, including said God □ therefore God is neither implied nor necessary, and may be shaved off and flushed by parsimony • there are viable alternatives to a definite earliest time, including an infinite past duration (which does not imply a contradiction however counter-intuitive), or a no-boundary, “edge-free” universe (which is not infinite in past duration) • the 1^st premise may be questionable or ambiguous in light of virtual particle pairs, quantum fluctuations, radioactive decay, spacetime foam/turbulence, the “pressure” of vacuum energy, the Casimir effect, Fomin’s quantum cosmogenesis (and successors), etc • God = not spatiotemporal incorporeal mind ⇒ such an “entity” would “be” no-where and no-when; thinking, decision making, acting, etc, would be impossible; such an entity could not be characterized as “free” • a supposed supernatural “beyond” is like an explanation that isn’t really much of an explanation to begin with (gaps, creative inventions) "If you want to believe in a personal creator God, based on your personal spiritual experiences, it's perfectly reasonable for you to do so, and you can ignore the arguments as to why there is no God, which are as flawed as the ones in favour." To understand why this is you should google up and read the mathematician David Hibert's famous thought experiment , "Hibert's Hotel". — John Gould Fishfry and, I expect most of the others that have responded to your thread, is/are perfectly familiar with Hilbert's Hotel, and why it is not an argument for any statement other than ' aren't infinities interesting ?' Ditto for Aristotle's notions of potential and actual infinities. Kalam has been discussed ad nauseam on this forum and its predecessor. You are very unlikely to come up with any arguments in its favour that have not already been considered and dismissed. Don't you think that, if it stood up to scrutiny, non-religious logicians might have noticed and written supportive papers about it? If you want to believe in a personal creator God, based on your personal spiritual experiences, it's perfectly reasonable for you to do so, and you can ignore the arguments as to why there is no God, which are as flawed as the ones in favour. But try not to be tempted into the hubris of believing that the correctness of your belief can be logically proven, and the non-believers are just too silly to see that. andrewk OptionsShare Your arguments dont seem more solid than the those made by a logical positivist. Why do you assert that your reason is infallible and can interpret the world as it really is by itself? That logical statements made by our fallible language can derrive at such mystical conclusions as proof for a personal creator? Why havent you ever even questioned the idea that Perhaps the whole idea of cause and effect is an illusion etc? Would you agree if someone Said that the law of cause and effect can only be applied to things that doesnt have independent and continuing being? I guess you would. Would you also agree that there then cant be any causal relationship between between entities that exist by their own power independent of the environment? John Gould What you are presuming is that actual infinities really exist. Actual infinity, however is merely an abstract notion , or, if you like, a fiction in the realm of the philosophy of mathematics which proposes that mathematical objects like , say the infinite sequence of negative numbers you refer to above can form a complete totality or "set", I.e. a given object that is a true actual infinity. Actual infinities, though, do not exist, they are not realities. To understand why this is you should google up and read the mathematician David Hibert's famous thought experiment , "Hibert's Hotel". What makes anyone thing the universe "began to exist?" For all we know it's like the negative integers ..., -3, -2, -1. Every element has a predecessor (or "cause" if you like) but there is no first What is the evidence for assumption P2? fishfry OptionsShare John Gould NB: with regard to why the creator of the universe was (a la the Kalam argument) a PERSoNAL God, I'll explain that for you in detail tomorrow when I have more time. And as for your vexed query re the notion of eternal punishment in Christian theology, please let's try to keep this discussion on track, an invitation to present and debate substantive philosophical/scientific objections to the Kalam cosmological argument. St Augustine's position on the question of "eternal punishment" is a totally separate issue, and at present I am not interested in explaining to you. It's something for a separate thread ( or your psychiatrist) Right? John Gould You say, "Why not say if the universe has a cause, that cause must be an effect of an effect of an effect" ? Here you are challenging img premise two ( P2) of the original Kalam argument : "The Universe BEGAN to exist"; you are asking why we should not suppose that the universe NEVER began to exist, rather, there was an infinite number of past cause -and -effect events prior to today? The reason is that an infinite number of things cannot exist. A POTENTIALLY infinite number of things can exist, but not an ACTUALLY infinite number of things. If one begins to argue than an ACTUALLY infinite number of things can exist certain absurdities inevitably result. The best illustration of this is "Hilbert's Hotel", the brainchild of the great German Mathematician David Hibert. Google "Hilbert's Hotel" for yourself, it a very accessible piece of reasoning written in clear and simple English ( even you Beebert will understand it and find the logic irrefutable)If you stil need further explanations/demonstrations re why ACTUAL infinities are fictions that exist only in the domain of mathematical discourse and not actual realities, I will provide them for you. Secondly, you ask why the universe could not have come into existence as a random phenomenon? Here you are challenging the first premise ( P1) of the traditional Kalam argument , namely : "Whatever begins to exist must have a cause for its beginning" Your objection is easy to rebut because to claim that something can just suddenly, randomly, "pop" into existence "ex nihilo" ,as it were, without any reasonable cause is worse than magic. When a magician pulls a rabbit out of a hat, at least you've got the hat, not to mention the magician. But if you deny P1. You are arguing that the whole universe just appeared at some point in the past for no reason whatsoever. But NOBODY - even you - SINCERELY believes that things, say a kangaroo or a Rolls Royce car or a sky-scraper just pop into being without a cause. Right ? "Saint Augustine put it -" Si Comprehendis non Deus est" - i.e. if you understand Him , then he is not God. " And yet this sadist called Augustine did Everything be could to "understand" this God with the help of a book, so that he could imagine him condemning as many People as possible to hell from before the foundation lf the world(before the foundation of the world... What an illogical contradiction in terms in this case) "P3. If the universe has a cause, then an uncaused, personal Creator of the universe exists, who sans the universe is beginningless, changeless, immaterial, timeless, spaceless and enormously Well this is a grotesque leap from the former statement in the formula that the universe has a cause. Why not as well say "If the universe has a cause, that cause must be an effect of an effect of an effect" or "If the universe has a cause, that cause must be random" or Whatever other stupid thing one might invent in one's stupid and proud head? I by the way claim that one can question that the universe even has a cause. And Why must this creator be personal? Why not say that he must be mindless and stupid and without any power because he was FORCED to create? Especially this is stupid. If he creates and is a cause, "he" becomes the the one who created. There are so many other things worthy to question "The Kalam argument provides strong evidence" This would probably even make Immanuel Kant laugh, despite his belief in synthetic judgements apriori or whatever. For one thing, to even be able to even claim that "Whatever begins to exist has a cause;" is true, one needs experience. But explain to me why this statement is even true? Also, if God is all those things Kalam said he is (timeless etc), does that mean then that from God's point of view, the world has eternally and timelessly existed? If not, isnt he then changing? John Gould My claim is that the Kalam argument is, to date, very strongly supported by philosophical reasoning and hard experimental, empirical scientific evidence and thus it gives us powerful grounds for believing in the EXISTENCE of a creator God. I did NOT say that the Kalam argument provides a means for us to KNOW anything about the nature of this God as He actually is in Himself. The Kalam argument is the product of human reasoning, it logically suggests that IN OUR CRUDE, VAGUE,LIMITED, AND INADEQUATE HUMAN TERMS we might describe the God ( divine being) who created the universe we inhabit 14 billion years ago as a beginningless, uncaused, atemporal, a -spatial, changeless, non-physical/immaterial, unconditional, enormously powerful, Personal creator. Insofar as you are asking WHO this God is -what is His true nature as He is in Himself that is an entirely different question; a question the Kalam argument does not seek to answer. All we can ultimately say is that this God -whom the Kalam argument gives us good reason to believe does exist is that with respect to humanity, He is transcendent and "wholly other". For us, He is utterly unknowable, unspeakable, unthinkable and totally incomprehensible; He Himself is forever hidden from our view. If He were not, then He would not be God; or, as Saint Augustine put it -" Si Comprehendis non Deus est" - i.e. if you understand Him , then he is not God. Briefly, there is, to be precise, no divine predicate/affirmation, no divine concept that contains in particular that which the God who created the universe is, there is only the divine subject as such and in Him the fullness of His divine affirmation. In short, The Kalam argument provides strong evidence for the existence of a divine creator God - God as a "known unknown". Metaphysician Undercover That seems unintelligible, for how can a first cause be simultaneously a final cause? — Brian A "Final cause" refers to a type of causation, the "final" does not refer to a temporal order. So there is nothing unintelligible about the first cause being a final cause. This would just be to say that the first cause is that type of cause, in comparison to an efficient cause for example. I dont like any arguments that obsesses itself with thinking in causes like that. I would rather call it one of the weakest arguments. And plus, what need do we have for this kind of pseudo-proofs? None. If one has Faith, one shall follow God. God doesnt reveal himself in "proofs" made by logical arguments. "P3. If the universe has a cause, then an uncaused, personal Creator of the universe exists, who sans the universe is beginningless, changeless, immaterial, timeless, spaceless and enormously Here is where Everything is lost. I would venture to say that the argument is even confused and weak right from the start. Part of learning to know God and realize he is seems to me to be to stop reasoning about him like that. Our common notion of causality requires the passage of time. What does it mean for something to "begin" to exist, or "have a cause" outside of time? — darthbarracuda Without time, the cosmological argument doesn't make sense. A cyclical universe could be posited to counter the Kalam argument. Each point in the cycle is both the end and the beginning. Since infinite regress doesn't come into the picture, we don't need to have a first cause. TheMadFool OptionsShare Our common notion of causality requires the passage of time. What does it mean for something to "begin" to exist, or "have a cause" outside of time? _db OptionsShare Brian A That seems unintelligible, for how can a first cause be simultaneously a final cause? Perhaps you mean that the final cause of the universe, viz. its teleology, demonstrates the hand of a Creator, as it were. But that would be the teleological argument, not the Kalam cosmological argument mentioned in the first post. The mere fact that the universe has a cause does not necessarily entail the view that the universe has a teleological purpose evidencing God. The latter may be so, but something beyond the poster's argument is necessary to establish it. Metaphysician Undercover As much as this does not demonstrate that God is "personal", there are ways to derive the conclusion that the first cause is very likely the type of cause which is commonly referred to as "final cause". Final cause is exemplified by freely willed actions. This is why it is often said that the existence of the universe is according to the will of God. Brian A Here is an objection that occurs to me: the Kalam cosmological argument does not necessarily lead to an omni-benevolent and personal God. Granting the premise that the universe must have a cause, what prevents us from holding the view that the first cause might be impersonal and indifferent (not omni-benevolent)? And if the the first cause is impersonal and indifferent, it does not follow that divine justice will be eventually rendered in reference to the actions of humanity. Therefore the usual view of God is not established. P1. Whatever begins to exist has a cause; P2. The universe began to exist; C1. Therefore, the universe has a cause. P3. If the universe has a cause, then an uncaused, personal Creator of the universe exists, who sans the universe is beginningless, changeless, immaterial, timeless, spaceless and enormously C2. Therefore, an uncaused, personal Creator of the universe exists, who sans the universe is beginningless, changeless, immaterial, timeless, spaceless and enormously powerful. One could take issue with any of the premises. Perhaps, in lieu of any actual constraints, spontaneous occurrence is possible. Perhaps the universe is beginningless. Perhaps the universe having a cause doesn't entail the existence of an uncaused, personal Creator who sans the universe is beginningless, changeless, immaterial, timeless, spaceless and enormously powerful. Michael OptionsShare John Gould I think the Kalam cosmological argument for the existence of God is one of the strongest defences for Theism that I have read. If anyone has any material objections to the Kalam proof , I would be interested in hearing them. What you are presuming is that actual infinities really exist. Actual infinity, however is merely an abstract notion , or, if you like, a fiction in the realm of the philosophy of mathematics which proposes that mathematical objects like , say the infinite sequence of negative numbers you refer to above can form a complete totality or "set", I.e. a given object that is a true actual infinity. Actual infinities, though, do not exist, they are not realities. To understand why this is you should google up and read the mathematician David Hibert's famous thought experiment , "Hibert's Hotel". — John Gould Thanks for your comments. I disagree with you on three points, summarized as follows. Even if I were using actual infinity, so what? After all, actual infinity is no weirder than an all-knowing, all-powerful, benevolent uncaused cause. 2. However, I am NOT using actual infinity. Your understanding of potential versus actual infinity is different than Aristotle's. My model "..., -3, -2, -1" only uses POTENTIAL infinity as defined by Aristotle. I will expound on this point in a moment. 3. Hilbert's Hotel (HH) doesn't apply. You are correct that HH assumes actual infinity. But my example only requires potential infinity. I don't need all the numbers (or rooms) to exist all at once. I only need that given one, I can identify the next. I never assume I have them all existing at once. That's potential infinity. Here is more detail, especially on point #2. 1. For the moment let me grant your (false) premise. Say I did need actual infinity (which I remind you I don't). So what? Craig wants us to conclude that there must be an uncaused cause, which he calls God. What if I call it absoulte infinity? I can't conceive of a worldview that would grant divinity but not infinity Cantor himself thought that his Absolute infinity was God. But Cantor's Absolute infinity is a lot bigger than the infinity of the natural numbers. 2. But #1 is irrelevant, since I don't use actual infinity, only potential. Let me outline the concept as seen by Aristotle. He said that we all have an intuition of the natural numbers 0, 1, 2, 3, 4, 5, ... Now the "dot dot dot" means something specific as defined by the Peano axioms Inductive axiom : Given a number n, there is a number n' called the "successor" of n. Another notion for the successor of n is n + 1. So if 0 exists, then 1 does. If 1 exists then 2 does. If 2 exists then 3 does. So if you want to know, does 43242342 exist? Then you can recursively drill down all the way back to 0, the base of the induction, and you can show that any particular number exists. There is never any claim that we have all of them together all at once . We can imagine they don't come into existence till we need them. All I need is n + 1 given n. If I need a million, I make a million. I never have them all at once. That is exactly Aristotle's definition of potential infinity . In the following quote, Aristotle is speaking of the endless regress 1/2, 1/4, 1/8, 1/16, etc. He says: "For the fact that the process of dividing never comes to an end ensures that this activity exists potentially, but not that the infinite exists separately." —Metaphysics, book 9, chapter 6. Now the sequence 1/2, 1/4, 1/8, ... may be identified with the the reverse sequence I gave earlier, ..., -3, -2, -1 by the simple mathematical trick of taking the base-2 logarithm of each fraction to get the corresponding negative integer. 1/2 maps to -1, and 1/4 maps to -2, and so forth. So these two examples are really the same example in different forms. Aristotle calls this infinity. Does 5 exist? Yes, if 4 exists. So we can prove that any number n exists. But we can't say that ALL the counting numbers exist all at the same time. That would be ACTUAL infinity. It was the genius of Cantor to take the huge conceptual leap and say, What if we allow actual infinity into math? That was his brilliant history-changing leap of the imagination. Cantor's insight was to write the following notation: {0, 1, 2, 3, ...} The braces symbolize the COMPLETED SET of natural numbers. The inductive axiom gives us 1 (given 0) then 2, then 3, and so forth. Axiom of Infinity says that there is a set containing all the natural numbers. That's actual infinity. We can summarize all this with a table. I apologize that I could not make the right column line up no matter what I did with tabs and spaces. Advice appreciated. Potential infinity Actual infinity Axiom of induction Axiom of Infinity Peano Cantor 0, 1, 2, 3, ... {0, 1, 2, 3, ...} n+1 given n All of them at once Negative integers Hilbert Hotel I hope this is helpful. 3. HH is just a popularized visualization of the fact that an infinite set may be placed into bijection with one of its proper subsets. In fact this can be taken as the definining property of infinite sets. You are right that HH does assume actual infinity. But we do NOT need actual infinity to define the natural numbers 1, 2, 3, ... in their usual order, or ..., -3, -2, -1 in their reverse order, or 1/ 2, 1/4, 1/8, ... if you use base 2 exponentials and logarithms. I don't need the strength of the axiom of infinity to give my example. Only the Peano axioms, which define potential infinity. Given n I need n + 1. I never need to complete the process. I only need to take the next step. So Hilbert's Hotel is not relevant here. fishfry OptionsShare John Gould Fish fry, Thank you for your detailed and very interesting response, though I do not want this thread to be diverted too far into the domain of the philosophy of mathematics; therefore let's put aside "Hibert's Hotel" and the notions of potential and actual infinities altogether shall we, because there is, in fact, no stipulation in the Kalam argument that the cause of the origin of the universe must be chronologically prior to that origin. Let's hypothesise instead, like Craig, that (for example) the Creator may be conceived causally, but not temporally prior to the origin of the universe such that the act of causing the universe to exist is SIMULTANEOUSLY with its beginning to exist? John Gould I agree with what you say about the Kalam argument and the "cutting edge of science", it is noteworthy, however, that in a sense, the history of 20th century cosmology can be seen as one failed attempt after another to avoid the absolute beginning of the universe predicted by the standard "Big Bang" model. That prediction is consistent with the Kalam argument and has now Stodden firm for nearly 100 years throughout a period of enormous change and tremendous advances in observational astronomy and creative theoretical work in astrophysics (?) You are right that the original Kalam cosmological argument refers only the space-time universe that we observe, that particular universe we human beings inhabit that is now believed by the majority of mainstream scientists to have come into being 14 billion or so years ago with the "Big Bang", etc. Regardless of what may exist beyond it, and regardless of any unanswered metaphysical questions about the true nature of absolute/ ultimate reality, our own universe is , in itself, no small, trifling thing , and in endeavouring to account for it the Kalam argument , in my opinion, is already addressing a formidable challenge. In short, limited in scope as it is merely to our universe, Kalam still, in my opinion, provides a vey strong, logically sound and rationally compelling case for the existence of a transcendent, supernatural Creator being. (And) that, for me, is quite enough food for thought just in itself !
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What is data structure? Definition, types, examples Every program relies on data and algorithms. Algorithms are sets of instructions that dictate how data is processed to produce meaningful results. Data, on the other hand, represents the information that algorithms operate on. Together, they form the foundation on which all software systems are built. Data structures play a crucial role in the interaction between algorithms and data. This article explores the fundamentals of data structures and how they function, showing the differences between various types with examples. Whether you're just starting out or looking to enhance your programming skills, this article will equip you with the knowledge needed to navigate the complexities of data structures. What is data structure? Data structure is a specialized format for organizing, processing, retrieving, updating, and storing data. Data structures serve as frameworks for arranging data for specific needs or objectives. They not only store the actual data values but also maintain information about how those values are related to each other. There are three main characteristics of a data structure that developers and data experts consider: correctness, time complexity, and space complexity. Characteristics of a data structure Correctness refers to the accuracy and reliability of the data structure's implementation. It should perform operations as expected, handle edge cases, and produce the right results. Time complexity is the efficiency of data structure operations in terms of the time required to perform them. It helps us understand how the runtime of an operation grows as the input size increases. Space complexity relates to the amount of memory an algorithm requires to perform its tasks efficiently on a particular data structure. Ideally, data structures should use as little memory as possible. This is particularly important in resource-constrained environments or when dealing with large datasets. Balancing these characteristics is essential for selecting data structures that meet the performance requirements and maintain accuracy and reliability in all scenarios. Data structure types This section will discuss the main data structure types, explaining their classification and functionalities. Check out the graph below to better understand how data structure types are categorized. Data structure classification Primitive data structures or types are the most basic data units available in programming languages. They include • integers: whole numbers—positive, negative, or zero, for example, 1, -30, 100, etc.; • floats: decimal numbers with a fractional part, like 3.14, -0.5, 7.0; • characters: individual characters, for example, “A,” “d,” and “1”; • strings: sequences of characters; and • booleans: binary values that can be either true or false, commonly used for conditional expressions. Primitive data structures represent simple values that cannot be broken down further into smaller components. They have a fixed size and format, and this predictability helps optimize memory usage. Primitive data types are also consistent across various programming languages, with only slight variations in naming or specific implementations. Non-primitive data structures are created with primitive data structures as their building blocks to efficiently organize and manage a collection of data. They can handle different data types and complex operations like searching, sorting, insertion, deletion, and more. Non-primitive data structures fall into two large categories: linear and non linear structures. Linear data structures vs. non linear data structures The following sections will explain each category in more detail. Linear data structures In a linear data structure, the elements are arranged sequentially in a particular order at one level. Each element has one predecessor and one successor, except for the first and last elements. This arrangement allows for a single uninterrupted run or iteration through the structure, starting from one end and progressing to the other. Linear data structures are typically straightforward and efficient for basic operations like adding, retrieving, or deleting. However, as the program becomes more complex, the limitations of this category become apparent. Even though linear data structures allow for single-run traversal, the process can still be complex as you have to visit each element one by one from the beginning. This results in a time complexity that increases linearly with the number of elements added. Let’s look at the types of linear data structures in more detail. Array data structure: Effective data storage and retrieval An array stores elements at contiguous memory locations—next to each other without gaps. It’s homogeneous, meaning that all data within a structure is of the same data type (only integers, characters, or others described earlier.) Other linear data structures can be homogeneous or heterogeneous depending on the programming language and their implementation. In an array, each element is associated with an index, a numerical identifier that indicates the element's position. Indexing begins at 0 for the first element and increments sequentially up to the array size minus one. For example, in an array of size 5, the indices range from 0 to 4, like in the picture below. Compared to other linear structures, a key advantage here is that accessing any element takes the same time regardless of its position or the array’s size. Depending on the programming language, arrays can be fixed or flexible in length. For example, in Java, you specify a constant size for an array, while in Python, you can create dynamic arrays. Arrays are commonly used in algorithms that require random access to elements, such as searching and sorting. They are also effective for storing lists of information (like dates or addresses), performing mathematical computations, and image processing. Additionally, arrays serve as the foundation for more complex data structures. Stack data structure: Managing undo and accessing recent elements A stack operates on the principle of Last In, First Out (LIFO), where inserting and retrieving data is possible from only one end. The last element added to the stack is the first one to be removed. Stack elements have two primary operations—push and pop. Pushing adds a new element, making it the top of the stack. Popping removes the topmost element first, exposing the next one as the new top of the stack. Regardless of the stack size, pushing or popping an element takes the same amount of time because there is always only one place to do the operation: the top of the stack. Stacks are particularly useful for managing data when the order of operations is important and when you need to access the most recently added elements first. A clear example is undo mechanisms in text editors, where each action performed by a user is recorded as a command and pushed onto the stack. When you initiate an undo operation, commands are popped from the stack, canceling the previous actions and restoring the document to its previous state. Queue data structure: Sequential processing of tasks or data A queue operates on the principle of First In, First Out (FIFO), where the first element added to the queue is the first one to be removed. In this case, unlike with a stack, inserting and retrieving data is done from different ends. In a queue, the elements are added at the back (enqueue) and removed from the front (dequeue). Similar to a stack, adding and removing an element takes the same time regardless of the queue size. Queues help process tasks or data in the order they were received. An example would be programs sending their print jobs, typically with only one printer available to process them sequentially. This printer handles each job in order, one at a time. Additionally, in event-driven programming, where events are processed as they occur, queues preserve the correct sequence of actions in the system. Linked list data structure: Flexibility with dynamically growing data A linked list consists of elements called nodes, each containing both data and a pointer (reference) to the next node in the sequence. The first node is called the head, and the last one has a null reference, indicating the end of the list. Each node can be placed at any available memory location, with the references between nodes enabling the traversal of the list. Linked list data structure Linked lists can grow or shrink in size dynamically, making them suitable for situations where the size of the data structure is not known in advance or may change. An example is a web browser history, where a webpage is represented as a node containing a reference to the next webpage visited. Non linear data structures Non linear data structures store elements not sequentially but rather in a hierarchical or arbitrary manner at different levels, often handling multiple connections. In such structures, you cannot traverse data in a single run. Non linear data structures are more difficult to implement but offer better performance for complex operations. They maintain constant time and space complexity as the size of the data grows. Non-linear structures are also more memory efficient as they allow the same connections to be shared among multiple elements. This reduces memory consumption compared to linear structures, where each element is connected to only one predecessor and/or successor. Let's dive into each type of non-linear data structure. Tree data structure: Depicting hierarchical relationships with categories A tree is a collection of elements (nodes) connected with edges that reflect parent-child relationships. Each node (except the root—the topmost node) has a single parent. Nodes without children are known as leaves. You’ve definitely seen this structure when looking at a family tree or organizational chart. The image above shows a binary tree where each node has at most two children. There are also trinary trees with three children per node and n-ary trees where each node can have up to n children, with n being a fixed positive integer. The height of a tree is the longest distance from the root to any leaf. The depth of a node is the length of the path from the root to that particular node. In other words, it is the level at which the node resides in the tree. Trees are commonly used in databases to represent categories and subcategories and in file systems—for hierarchically organizing files and directories. Heap data structure: Prioritizing tasks by importance A heap is a binary tree where the topmost node always has the highest or the lowest value in the whole tree. This property is called the heap property or heap invariant. There are two types of heaps. In a min-heap, the value of each node is less than or equal to the values of its children. In a max-heap, the value of each node is greater than or equal to the values of its children. With the heap data structure, we can implement priority queues, where elements have different importance. For example, job scheduling applications need to organize tasks based on urgency or deadline, and that’s where heaps come in handy. Trie data structure: Efficient storage and retrieval of text strings A trie, also referred to as a prefix tree or a keyword tree, is used for storing and retrieving a set of strings. In this structure, each node contains a character—typically, a letter of the A prefix is one or more characters that form the beginning part of one or more strings. For example, in the image above, b is a common prefix for the words ball and bee, while an is a prefix for ant and angel. Tries are applied in cases where efficient storage and retrieval of strings are required, particularly when dealing with large datasets or dictionaries. A common example is autocompletion, where the trie structure allows quick suggestions based on partial input strings. In spell-checking algorithms, tries can also determine whether a given string is a valid word and can suggest corrections. Graph data structure: Visualizing relationships in complex networks A graph is a data structure used to represent relationships between entities. It consists of a set of nodes, also known as vertices, and each vertex connects to others through edges. One of the fundamental distinctions between graphs and trees is that graphs can contain cycles, while trees cannot. A cycle is a path in the graph that starts and ends at the same vertex, traversing edges without repeating any vertex. The graph structure can depict complex network relationships, enabling efficient analysis and processing of interconnected data. For example, graphs can represent computer networks, molecular structures, and circuit layouts. Hash table data structure: Streamlining searches and sorting operations A hash table or hash map stores key-value pairs in an array. The key is a unique identifier to access or retrieve the value, which is the associated data or information. Hash table data structure The hash function takes a key as input and produces a unique numerical value called a hash code that serves as an index for a particular data element. This mapping process is deterministic, meaning the same key will always be hashed to the same index, facilitating rapid lookup and retrieval of values. We use hash tables to quickly map keys to specific locations in the array. The goal of the hash function is to distribute keys appropriately across the array indices, minimizing collisions (situations where multiple keys map to the same index). Hash tables are common in scenarios where efficient storage and retrieval of key-value pairs are required — for example, caching systems. Additionally, we employ hash tables in data processing tasks such as indexing and deduplication. The evolution of data structures Linear data structures emerged to address simple storage and retrieval needs, while non-linear structures became necessary to represent complex relationships and optimize algorithmic performance. With the increasing volumes of data being generated and evolving domain-specific challenges, there will be a continued focus on developing data structures that offer improved efficiency, scalability, and resilience. Future advancements in areas like artificial intelligence, Internet of Things (IoT), blockchain, and edge computing may lead to the development of specialized data structures that cater to the unique requirements of these domains. It might involve advancements in data structure design, algorithmic optimizations, and parallel processing techniques. Adaptability and openness to change are essential in the ever-changing field of technology. Overlooking this may result in performance bottlenecks, limited problem-solving abilities, and suboptimal solutions that are slower, consume more memory, and have poorer scalability.
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(All citation calculations based on Total Publications Data From 483 Publications This is the number of all publications for the study. Total Citations Data From 425 Publications This is the number of citations to all publications for the project. Mean Citations Per Publication Data From 425 Publications The total number of citations divided by the total number of publications. Age-weighted Mean Citations Per Publication Data From 483 Publications Age-weighted Mean Citations Per Publication. Median Citations Data From 425 Publications The median number of citations for the project. Data From 425 Publications h-index is the largest number h such that h publications from a study have at least h citations. Data From 425 Publications The number of publications that have at least 100 citations. Data from 483 publications. (What does this mean?) Data from 483 publications. (What does this mean?) Data from 425 publications. (What does this mean?) Data from 425 publications. (What does this mean?)
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How do you calculate shear rate and viscosity? The shear rate is the velocity of the upper plate (in meters per second) divided by the distance between the two plates (in meters). Its unit is [1/s] or reciprocal second [s-1]. According to Newton’s Law, shear stress is viscosity times shear rate. Therefore, the viscosity (eta) is shear stress divided by shear rate. How do you calculate the shear rate of a viscometer? When we perform rheological experiments with a viscometer, usually it gives us shear stress values corresponding to particular revolution per minute (rpm). Then we convert rpm into a shear rate by multiplying a factor 1.7. How do you calculate extruder shear rate? To calculate the shear rates between the screw flight and the barrel wall, use Eqn (8.1) with h = distance between the screw flight and the barrel wall. where τ = Shear stress = F/A = Force applied per unit area. An example of the resistance factor in a rod die on the flow rate is shown in Figure 8.1. What is shear viscosity? 1. A coefficient that characterizes the viscous properties of a fluid and is related to the absorption (loss) of energy (or else, damping) due to the presence of velocity gradients in the fluid. What is shear measured in? Physical quantities of shear stress are measured in force divided by area. In SI, the unit is the pascal (Pa) or newtons per square meter. In United States customary units, shear stress is also commonly measured in pounds-force per square inch or kilopounds-force per square inch. What are the units of shear rate? Shear rate is the speed of deformation in the shear mode (which is typical of fluids and can be represented as layers sliding one onto another). It is expressed (SI units) in reciprocal seconds [1/ s], since it derives from radiant per second, with radiant being dimensionless (just a number). Does viscosity depend on shear rate? The viscosity remains constant with changing shear rate. However, most fluids are non-Newtonian, that is, their viscosities are a function of shear rate. Therefore, changing the geometry, such as the distance between the two plates, above, and/or the speed of flow, may significantly change the viscosity. What is the unit of shear rate? The SI unit of measurement for shear rate is s −1, expressed as “reciprocal seconds” or “inverse seconds”. The shear rate at the inner wall of a Newtonian fluid flowing within a pipe is. where: .γ is the shear rate, measured in reciprocal seconds; v is the linear fluid velocity; d is the inside diameter of the pipe. How is shear strain calculated? Shear Strain Example First, measure the original length. Measure the original length, often considered the height of a square object. Next, measure the deformation. After a shear force has been applied to the object, measure the deformation. Finally, calculate the shear strain. Calculate the shear strain by dividing the deformation by the original length. What is shear ratio? Shear modulus, in materials science, is defined as the ratio of shear stress to shear strain. The shear modulus value is always a positive number and is expressed as an amount of force per unit area. What is wall shear stress and shear strain rate? wall shear stress is the force the ‘tear’ the fluid near the wall due to no slip boundary condition. shear strain rate is the rate of fluid being ‘teared’ . higher shear strain rate means the fluid is being ‘teared’ more compare to lower shear strain rate when the time is constant.
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Excel Formula Python - Select Cell with Sum of Values In this tutorial, we will learn how to write an Excel formula in Python that selects a cell from a range based on the sum of values. This formula can be useful when you have a range of cells and you want to find the cell that has a cumulative sum of values up to a specific number. To achieve this, we will use the INDEX and MATCH functions in Excel, which can be implemented in Python using the openpyxl library. The formula we will use is as follows: =INDEX(F12:F19, MATCH(TRUE, MMULT(--(F12:F19<=48958000), TRANSPOSE(COLUMN(F12:F19)^0))=1, 0)) Let's break down the formula step by step: 1. The MMULT function is used to calculate the cumulative sum of values in the range F12:F19. It converts the values to an array of 1s and 0s, where 1 represents the cells that have a cumulative sum less than or equal to the specified number. 2. The TRANSPOSE function is used to transpose the array of 1s and 0s, so that it can be multiplied with another array. 3. The COLUMN function is used to create an array of numbers from 1 to the number of cells in the range F12:F19. 4. The ^0 operation is used to convert the array of numbers to an array of 1s. 5. The -- operation is used to convert the array of 1s and 0s to an array of TRUE and FALSE values. 6. The MMULT function multiplies the array of TRUE and FALSE values with the transposed array of 1s and 0s, resulting in an array of 1s and 0s where 1 represents the cells that have a cumulative sum less than or equal to the specified number. 7. The MATCH function is used to find the position of the first occurrence of TRUE in the array of 1s and 0s. 8. The INDEX function is used to return the value from the range F12:F19 at the position found by the MATCH function. To use this formula in Python, you can install the openpyxl library and use the following code: import openpyxl workbook = openpyxl.load_workbook('path/to/excel/file.xlsx') worksheet = workbook['Sheet1'] cell_value = worksheet['F12'].value Replace 'path/to/excel/file.xlsx' with the actual path to your Excel file, and 'Sheet1' with the name of your worksheet. The formula will be applied to the range F12:F19, and the selected cell value will be printed. In conclusion, this tutorial has shown you how to write an Excel formula in Python to select a cell from a range based on the sum of values. This can be useful in various scenarios where you need to find a specific cell that meets certain criteria. By using the openpyxl library, you can easily implement this formula in your Python code and work with Excel files programmatically. An Excel formula =INDEX(F12:F19, MATCH(TRUE, MMULT(--(F12:F19<=48958000), TRANSPOSE(COLUMN(F12:F19)^0))=1, 0)) Formula Explanation This formula uses the INDEX and MATCH functions to select the cell from F12 to F19 that has a sum of values up to 48,958,000. Step-by-step explanation 1. The MMULT function is used to calculate the cumulative sum of values in the range F12:F19. It converts the values to an array of 1s and 0s, where 1 represents the cells that have a cumulative sum less than or equal to 48,958,000. 2. The TRANSPOSE function is used to transpose the array of 1s and 0s, so that it can be multiplied with another array. 3. The COLUMN function is used to create an array of numbers from 1 to the number of cells in the range F12:F19. 4. The ^0 operation is used to convert the array of numbers to an array of 1s. 5. The -- operation is used to convert the array of 1s and 0s to an array of TRUE and FALSE values. 6. The MMULT function multiplies the array of TRUE and FALSE values with the transposed array of 1s and 0s, resulting in an array of 1s and 0s where 1 represents the cells that have a cumulative sum less than or equal to 48,958,000. 7. The MATCH function is used to find the position of the first occurrence of TRUE in the array of 1s and 0s. 8. The INDEX function is used to return the value from the range F12:F19 at the position found by the MATCH function. For example, if we have the following data in cells F12 to F19: | F | | | | 100 | | 200 | | 300 | | 400 | | 500 | | 600 | | 700 | | 800 | The formula =INDEX(F12:F19, MATCH(TRUE, MMULT(--(F12:F19<=48958000), TRANSPOSE(COLUMN(F12:F19)^0))=1, 0)) would return the value 400, which is the cell from F12 to F19 that has a sum of values up to
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Multiple Charts In One Excel 2024 - Multiplication Chart Printable Multiple Charts In One Excel Multiple Charts In One Excel – You may create a multiplication graph in Excel by using a format. You will discover a number of types of web templates and learn to file format your multiplication graph using them. Here are some tricks and tips to produce a multiplication graph. Once you have a design, all you have to do is version the formula and mixture it inside a new cell. You may then take advantage of this formula to grow some numbers by yet another set. Multiple Charts In One Excel. Multiplication dinner table web template You may want to learn how to write a simple formula if you are in the need to create a multiplication table. Initially, you have to locking mechanism row one of the header column, then grow the telephone number on row A by mobile B. A different way to develop a multiplication kitchen table is to apply mixed personal references. In this instance, you would probably enter in $A2 into line A and B$1 into row B. The outcome is a multiplication desk having a formula that really works both for columns and rows. If you are using an Excel program, you can use the multiplication table template to create your table. Just open the spreadsheet with the multiplication desk template and change the label for the student’s label. You can also adjust the sheet to fit your personal needs. There is an solution to affect the shade of the cells to change the look of the multiplication desk, way too. Then, you can change all the different multiples for your needs. Creating a multiplication graph or chart in Stand out When you’re utilizing multiplication desk software program, it is possible to create a easy multiplication kitchen table in Stand out. Merely build a page with columns and rows numbered in one to 35. Where rows and columns intersect may be the respond to. For example, if a row has a digit of three, and a column has a digit of five, then the answer is three times five. The same goes for the Initially, you may enter into the amounts that you have to flourish. For example, if you need to multiply two digits by three, you can type a formula for each number in cell A1. To create the figures larger sized, pick the cellular material at A1 and A8, after which click on the proper arrow to decide on a variety of cellular material. You can then type the multiplication solution within the tissue inside the other rows and columns. Gallery of Multiple Charts In One Excel How To Plot Multiple Data Sets In One Excel Chart YouTube How To Quickly Make Multiple Charts In Excel YouTube Multiple Bar Charts On One Axis In Excel Super User Leave a Comment
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: .. : : 70 : : 2017 : .. // . 70. .: , 2017. .136-167. URL: https://doi.org/10.25728/ubs.2017.70.6 , , , cellular automaton, autonomous agents, reflexive agents, agents' formation control . , , , . -- . The article deals with the algorithm for a distributed organization of the agents' formation defined by a graph and the numerical simulation of such algorithm. Agents move through terrain with many random obstacles (``random landscape''). At first, we describe the continuous statement of the two-criteria minimization problem. The first criterion is the agent's route time. The second criterion is the closeness of agents' formation to the desired one. Next, we introduce a cellular automaton simulating the movement of agents for obtaining quasi-optimal solutions of the problem. The cellular automaton has one-dimensional and two-dimensional representations. Agents use reflexion to predict the motion of other agents. Then we compare obtained solutions with optimal ones for different types of random landscapes via numerical experiment. At this point, we obtain the empirical distribution for the time of an agent's exit to finish point. Finally, we find the relation between a type of random landscape through which agents move and the quality of agents' formation. PDF - : 3073, : 2815, : 23.
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What is 98 Celsius to Fahrenheit? - ConvertTemperatureintoCelsius.info 98 Celsius to Fahrenheit Conversion: Understanding the Relationship Between Celsius and Fahrenheit When it comes to measurements of temperature, Celsius and Fahrenheit are two commonly used units. Each unit provides a different scale for measuring temperature, and the conversion between them can be confusing at times. In this article, we will explore the relationship between Celsius and Fahrenheit and specifically answer the question, “What is 98 Celsius to Fahrenheit?” Understanding Celsius and Fahrenheit First, let’s understand the basics of Celsius and Fahrenheit. Celsius is a metric unit for measuring temperature and is commonly used worldwide. It is based on the freezing point of water at 0 degrees and the boiling point of water at 100 degrees. On the other hand, Fahrenheit is a measurement unit commonly used in the United States. It is based on a scale in which the freezing point of water is 32 degrees and the boiling point is 212 degrees. The Conversion Formula To convert Celsius to Fahrenheit, you can use the following formula: °F = (°C × 9/5) + 32 This formula allows you to easily convert temperatures from Celsius to Fahrenheit and vice versa. Now, let’s apply this formula to the specific question at hand: “What is 98 Celsius to Fahrenheit?” Converting 98 Celsius to Fahrenheit Using the conversion formula mentioned earlier, we can calculate the equivalent temperature in Fahrenheit for 98 degrees Celsius: °F = (98 × 9/5) + 32 °F = (176.4) + 32 °F = 208.4 Therefore, 98 degrees Celsius is equivalent to 208.4 degrees Fahrenheit. This simple calculation provides the answer to the initial question and demonstrates the conversion between Celsius and When using this information in everyday life, it’s important to remember that understanding temperature measurements in both Celsius and Fahrenheit can be beneficial. Whether you’re traveling to a different country or simply using a recipe from another part of the world, having a grasp of both temperature units can be valuable. Additionally, with the increasing globalization and interconnectedness of the world, having this knowledge can foster better communication and understanding across different cultures and regions. It can also aid in scientific and academic pursuits, where knowledge of temperature measurements is essential. In conclusion, understanding the conversion between Celsius and Fahrenheit is a valuable skill that can be applied in various aspects of life. By following the simple conversion formula and learning how to convert temperatures from one unit to another, individuals can expand their knowledge and adapt to different measurement systems. So, the next time you come across the question, “What is 98 Celsius to Fahrenheit?” you’ll be equipped with the knowledge to provide a definitive answer.
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ES386-15 Dynamics of Vibrating Systems Introductory description ES386-15 Dynamics of Vibrating Systems Module aims Vibrations exert a significant influence on the performance of the majority of engineering systems. All engineers should understand the basic concepts and all mechanical engineers should be familiar with the analytical techniques for the modelling and quantitative prediction of behaviour. Thus, this module provides students with fundamental skills necessary for the analysis of the dynamics of mechanical systems, as well as providing opportunities to apply these skills to the modelling and analysis of vibration. This third-year module is mandatory for students pursuing a degree in Mechanical Engineering, building upon competences acquired earlier in the course. This module introduces students to the use of Lagrange’s equations (applied to 1D and 2D systems only for this module) and to techniques for modelling both lumped and continuous vibrating systems. It includes some coverage of approximate methods both as an aid to physical understanding of the principles and because of their continuing usefulness. The module assumes basic understanding of mechanics and linear algebra consistent with the level of Year 2 modules. At the end of the module students should have a sound understanding of the wide application of vibration theory and of the underlying physical principles. In particular, they should be able to use either Newtonian or Lagrangian mechanics to analyse 2D systems, and to determine the response of simple damped and undamped multi-degrees of freedom (DOF) systems to both periodic and aperiodic excitation. They should also be familiar with engineering solutions for measuring and influencing vibrational behaviour. Outline syllabus This is an indicative module outline only to give an indication of the sort of topics that may be covered. Actual sessions held may differ. • Generalised co-ordinates, Lagrange's equation (including preliminary study of other classical methods) • General application of the Lagrange equation to vibrating systems • Multi-degree of freedom systems: lumped system models, continuous system models; geared and branched systems; reduction of an n-DOF system to a set of n single-DOF systems; principal co-ordinates • Matrix methods of analysis: conservative and non-conservative (damped) systems; determination of principal co-ordinates • Modelling of damping: hysteretic, Coulomb, viscous; measurement of damping factor • Forced vibration: harmonic excitation of multi-DOF systems; shaft whirling; transmissibility; vibration isolation; non-harmonic and arbitrary excitation (convolution integral) • Approximate methods e.g. Rayleigh's method, Dunkerley's method Learning outcomes By the end of the module, students should be able to: • 1. Model planar mechanical systems using Newton’s and Lagrange’s equations: Determine appropriate co-ordinate systems, analyse vibrations. • 2. Abstract more complex engineering mechanisms: analyse using lumped system models or simple distributed mass and stiffness models. Use and justify standard methods and approximations for extended and continuous vibrating systems. • 3. Evaluate the natural frequencies and modes of vibration of a multi-degree of freedom linear system. • 4. Determine and analyse the free and forced response of single-degree of freedom systems to periodic and aperiodic excitation, as well as the effects of linear and non-linear damping on the system behaviour. • 5. Evaluate complex (multi-degree of freedom) undamped or damped systems numerically, using a systematic approach to analyse the natural frequencies and modes, and the response of the system to periodic and aperiodic excitations. • 6. Demonstrate a sound understanding of the application of vibration analysis to key engineering systems. Indicative reading list 1. Theory of Vibration with Applications, by W. T. Thomson and M. D. Dahleh. Publisher: Pearson. Fifth edition, 1998. ISBN-10: 013651068X, ISBN-13: 9780136510680. 2. Principles of Vibration, by B. H. Tongue. Publisher: Oxford University Press. Second edition, 2002. ISBN-10: 0195142462. 3. Engineering Vibrations, by D. J. Inman. Publisher: Pearson. Fourth international edition, 2013. ISBN-10: 0273768441, ISBN-13: 9780273768449. 4. Mechanical vibrations, by S. S. Rao, Fook Fah Yap. Publisher: Prentice Hall. Fifth edition in SI units, 2011. ISBN-10: 9810687125, ISBN-13: 9789810687120 5. Vibrations, by B. Balachandran, E. B. Magrab. Publisher: Cengage. Second international SI edition, 2009. ISBN10: 0495411256, ISBN-13: 9780495411253. View reading list on Talis Aspire Subject specific skills SSS4: Ability to apply relevant practical and laboratory skills. SSS8: Ability to be pragmatic, taking a systematic approach and the logical and practical steps necessary for, often complex, concepts to become reality. Transferable skills TS1: Numeracy: apply mathematical and computational methods to communicate parameters, model and optimize solutions. TS2: Apply problem solving skills, information retrieval, and the effective use of general IT facilities. TS3: Communicate (written and oral; to technical and non-technical audiences) and work with others. TS7: Overcome difficulties by employing skills, knowledge and understanding in a flexible manner. Study time Type Required Lectures 30 sessions of 1 hour (25%) Seminars 2 sessions of 1 hour (2%) Practical classes 1 session of 2 hours (2%) Private study 86 hours (72%) Total 120 hours Private study description Guided independent learning, assignment preparation, etc 86 hours. No further costs have been identified for this module. You must pass all assessment components to pass the module. Students can register for this module without taking any assessment. Assessment group D4 Weighting Study time Eligible for self-certification Vibration computational and analysis assignment 30% Yes (extension) Matlab code submitted on Matlab Grader (10% of module credit) and a brief computational report (1000 words, 20% of module credit) Online Examination 70% No QMP test (2x 1h) ~Platforms - AEP,QMP • Online examination: No Answerbook required • Students may use a calculator • Engineering Data Book 8th Edition Feedback on assessment • Feedback during laboratory sessions • Feedback on assignments. • Model solutions to exam type questions. • Support through advice and feedback hours. • Cohort level feedback on examinations To take this module, you must have passed: This module is Core for: • Year 3 of UESA-H310 BEng Mechanical Engineering • Year 3 of UESA-H315 BEng Mechanical Engineering • Year 4 of UESA-H314 BEng Mechanical Engineering with Intercalated Year • Year 3 of UESA-HH35 BEng Systems Engineering • Year 3 of UESA-HH36 BEng Systems Engineering • Year 4 of UESA-HH34 BEng Systems Engineering with Intercalated Year • Year 3 of UESA-H311 MEng Mechanical Engineering • UESA-H316 MEng Mechanical Engineering □ Year 3 of H315 Mechanical Engineering BEng □ Year 3 of H316 Mechanical Engineering MEng • Year 4 of UESA-H317 MEng Mechanical Engineering with Intercalated Year • UESA-HH31 MEng Systems Engineering □ Year 3 of HH31 Systems Engineering □ Year 3 of HH35 Systems Engineering • Year 4 of UESA-HH32 MEng Systems Engineering with Intercalated Year This module is Core optional for: • Year 3 of UESA-H115 MEng Engineering with Intercalated Year • UESA-H317 MEng Mechanical Engineering with Intercalated Year □ Year 3 of H317 Mechanical Engineering with Intercalated Year □ Year 4 of H317 Mechanical Engineering with Intercalated Year • Year 4 of UESA-HH32 MEng Systems Engineering with Intercalated Year • Year 3 of UESA-H11L Undergradaute Engineering (with Intercalated Year) This module is Optional for: • Year 3 of UESA-H113 BEng Engineering • Year 3 of UESA-H114 MEng Engineering • Year 4 of UESA-H115 MEng Engineering with Intercalated Year • UESA-H11L Undergradaute Engineering (with Intercalated Year) □ Year 3 of H11L Engineering (with Intercalated Year) □ Year 4 of H11L Engineering (with Intercalated Year) This module is Option list A for: • Year 4 of UESA-H111 BEng Engineering with Intercalated Year • Year 3 of UESA-H112 BSc Engineering
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Stock y has a beta of 13 and an expected return of 185, Financial Management Stock Y has a beta of 1.3 and an expected return of 18.5%. Stock Z has a beta of 0.70 and an expected return of 12.1%. If the risk-free rate is 8% and the market risk premium is 7.5%, are these stocks correctly priced? If not, what would the risk-free rate have to be for the two stocks to be correctly priced? Request for Solution File Ask an Expert for Answer!! Financial Management: Stock y has a beta of 13 and an expected return of 185 Reference No:- TGS01039321 Expected delivery within 24 Hours
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Frontiers | The value of generalized linear mixed models for data analysis in the plant sciences • ^1Department of Plant Pathology, The Ohio State University, Wooster, OH, United States • ^2Center for Integrated Fungal Research, Department of Entomology and Plant Pathology, North Carolina State University, Raleigh, NC, United States Modern data analysis typically involves the fitting of a statistical model to data, which includes estimating the model parameters and their precision (standard errors) and testing hypotheses based on the parameter estimates. Linear mixed models (LMMs) fitted through likelihood methods have been the foundation for data analysis for well over a quarter of a century. These models allow the researcher to simultaneously consider fixed (e.g., treatment) and random (e.g., block and location) effects on the response variables and account for the correlation of observations, when it is assumed that the response variable has a normal distribution. Analysis of variance (ANOVA), which was developed about a century ago, can be considered a special case of the use of an LMM. A wide diversity of experimental and treatment designs, as well as correlations of the response variable, can be handled using these types of models. Many response variables are not normally distributed, of course, such as discrete variables that may or may not be expressed as a percentage (e.g., counts of insects or diseased plants) and continuous variables with asymmetrical distributions (e.g., survival time). As expansions of LMMs, generalized linear mixed models (GLMMs) can be used to analyze the data arising from several non-normal statistical distributions, including the discrete binomial, Poisson, and negative binomial, as well as the continuous gamma and beta. A GLMM allows the data analyst to better match the model to the data rather than to force the data to match a specific model. The increase in computer memory and processing speed, together with the development of user-friendly software and the progress in statistical theory and methodology, has made it practical for non-statisticians to use GLMMs since the late 2000s. The switch from LMMs to GLMMs is deceptive, however, as there are several major issues that must be thought about or judged when using a GLMM, which are mostly resolved for routine analyses with LMMs. These include the consideration of conditional versus marginal distributions and means, overdispersion (for discrete data), the model-fitting method [e.g., maximum likelihood (integral approximation), restricted pseudo-likelihood, and quasi-likelihood], and the choice of link function to relate the mean to the fixed and random effects. The issues are explained conceptually with different model formulations and subsequently with an example involving the percentage of diseased plants in a field study with wheat, as well as with simulated data, starting with a LMM and transitioning to a GLMM. A brief synopsis of the published GLMM-based analyses in the plant agricultural literature is presented to give readers a sense of the range of applications of this approach to data analysis. 1 Introduction Whether one is conducting a planned experiment with replication, blocking, and randomization or collecting observational data such as from a survey, data analysis needs to be carried out in order to reach conclusions (Schabenberger and Pierce, 2002). We take it as axiomatic that appropriate statistical methods are required to interpret the collected data from any study. Of course, “appropriate” can mean many different things, depending on the situation and context of interest. Generalized linear mixed models (GLMMs) are becoming quite popular for data analysis in agriculture and other disciplines (Gbur et al., 2012; Ruíz et al., 2023); however, there is a lot of uncertainty about which GLMM to use, how to fit GLMMs to data, and how to interpret the obtained results. Our objectives here were to provide some background on the statistical modeling of data, particularly for data collected in the agricultural or the plant sciences, to explain GLMMs with examples, and to discuss several important issues when using GLMMs that may be initially not appreciated by the data analyst. Several other references are valuable for learning more about GLMMs and for conducting analyses of a wide range of datasets from different experimental designs (Littell et al., 2006; Bolker et al., 2009; Zuur et al., 2009; Gbur et al., 2012; Stroup, 2013; Brown and Prescott, 2015; Stroup et al., 2018; Gianinetti, 2020; Li et al., 2023; Ruíz et al., 2015, 2023). For those who wish to learn more theory, as well as applications, Stroup (2013) is an indispensable reference. Molenberghs and Verbeke (2010) and McCulloch and Searle (2001) provided considerably more theory about mixed models, but may be of less value to the reader of this article. This article is heavily influenced by Stroup (2013, 2015) and the material in Chapters 11–13 in Stroup et al. (2018). We focus on an example involving the percentage of plants infected by a particular disease in a field study, although the methods can be applied to any discrete data where percentages can be calculated. Gianinetti (2020) and Li et al. (2023) are other excellent references for the GLMM-based analysis of data expressed as percentages. Readers should refer to Gbur et al. (2012); Stroup (2013), and Ruíz et al. (2023) for details on the analysis of continuous data using GLMMs. For those interested in broader issues related to statistical analysis in horticulture, including commonly made errors, we recommend Kramer et al. (2016). Below, we start with some historical background on data analysis, followed by a discussion on non-normal statistical distributions and then presentations on models for data with normal and non-normal distributions. Model expansions and alternatives are given for select experimental designs. We place an emphasis on the different methods of model fitting and the interpretation of the estimated parameters for GLMMs, demonstrated with an example dataset. Some major challenges in the use of GLMMs are presented. Before the conclusions, a list of cases in the literature where GLMMs were used is 2 Background: from ANOVA to GLMMs Analysis of variance (ANOVA), including the special case of t-tests, and linear regression have been the foundations for data analysis in agriculture and other disciplines for over a century (Fisher, 1918, 1935; Cochran and Cox, 1957; Steel and Torrie, 1960). The pioneering statistical works of Fisher, Yates, Cochran, and Snedecor, among others, coupled with the advances in statistical software and the increased speed and memory of computers, have given researchers a large toolbox of methods to describe data, predict outcomes, and make inferences about hypotheses of interest (Schabenberger and Pierce, 2002; Gbur et al., 2012). ANOVA can be considered a special case of multiple linear regression (Speed, 2010), where several predictor or explanatory variables (X[1], X[2], etc.) are binary (0, 1). This allows data analysis to be couched in terms of linear modeling of the response variables as functions of predictor variables, an approach that still dominates today. An explanatory variable can be discrete, generally known as a classification or class variable, or simply a factor, and consists of two or more distinct levels (e.g., cultivar 1 and cultivar 2 or treatment A and treatment B). Otherwise, the explanatory variable can be continuous (e.g., temperature), often called a covariate or a covariable. Continuous variables can be treated as discrete variables in a statistical model, depending on the objective and the manner in which the experiment was conducted. Models consist of fixed-effects and/or random-effects variables (Schabenberger and Pierce, 2002). For the fixed-effects variables, the levels (categories or groups) in the study represent all possible levels of the factor or all the levels of interest by the investigator (i.e., they were selected for study because they are of specific interest). Examples would be fungicide treatment, biocontrol treatment, pathogen inoculum dose, temperature, cultivar, or the nitrogen level in a fertilizer. For the random-effects variables, the levels in the study represent only a random sample of a larger set of levels (i.e., a sample from a distribution of effects). Examples could include the location (environment), a block, or a plot in a field study. A given variable could be considered fixed or random, depending on the circumstances. For instance, plant genotype could be considered a random-effects variable if a sample of a population of genotypes was randomly selected for study, with the goal of characterizing the mean (expected value) and the variability of the response variable (e.g., yield); on the other hand, genotype could be considered as a fixed-effects variable if the investigator was strictly interested in the responses of those particular genotypes. Similar arguments can be made about location–year (“environment”) as being fixed or random (see Chapter 6 in Littell et al., 2006). The term mixed model is used when there are both fixed- and random-effects variables in the model. More specifically, these are known as linear mixed models (LMMs) when the response variable (e.g., yield, biomass, or disease severity) is considered to be normally distributed and the mean (expected value) is modeled directly as a function of the fixed- and random-effects terms (Stroup, 2013). This is explained in the next section, which includes a more technical description of LMMs. A special case of a LMM is the linear model (LM), in which there are no random effects except for the residual; another special case is the random-effects model, in which there are no fixed-effects variables. The concept of random effects goes back to Fisher (1918, 1935), possibly earlier (but less formally), with subsequent important early contributions by Yates (1940) and Eisenhart (1947), as well as many others. For many years, the standard approach to handling random effects in standard software (or by brute force work on a calculator in the very early days) was to fit a LM to the data using ordinary least squares as if all effects were fixed and then perform post-model-fitting calculations to estimate the mean squares and variances for the random-effects variables (Steel and Torrie, 1960; Littell et al., 2006). The latter are used (automatically) to estimate the standard errors (SEs) of the least squares means and the mean differences for the fixed-effects terms and also to test for factor effects. This approach works fine for many special cases (e.g., balanced randomized complete block designs); however, even for a design such as the popular split plot or split plot with blocks, certain SEs cannot be calculated correctly (Littell et al., 2006). For incomplete block designs (i.e., where each block does not contain all the treatments), recovering all of the available information in a study (e.g., treatment effects within and between blocks) requires some tedious post-model-fitting calculations when all factors (including blocks) are considered fixed in the model (Yates, 1940). Many other situations often cannot be analyzed satisfactorily with the mean square, post-model-fitting adjustment approach, at least not without major approximations. Examples include repeated measures in time and space or any situation when complex correlations of data need to be specified or estimated. The methodology of treating random effects fully as random in the model-fitting process was developed by Henderson in a series of papers (e.g., Henderson, 1950, 1953, 1984), with extensive theory by Harville (1977) and others, especially Laird and Ware (1982). The model-fitting approach is likelihood-based (rather than least squares/mean square-based), with an assumed normal distribution for the response variable conditional on the random effects (see below). Except for special (simple) cases, model fitting is iterative, facilitated by fast computers with large memories to carry out the multiple iterations in a rapid manner. It was not until the mid-1990s that (relatively) easy-to-use software was developed to fit LMMs to data based on the likelihood principle. The MIXED procedure in SAS is especially important here (Littell et al., 2006). Presently, several R packages, such as “lme4” and “nlme,” can be used (Galecki and Burzykowski, 2013; Bolker et al., 2022). GenStat (VSN International, Hemel Hempstead, UK) has been able to fit LMMs using the likelihood methodology for many years. These programs have opened up a great diversity of applications in data analysis across all disciplines that go far beyond what was possible with the traditional ANOVA approach in the older software (Brown and Prescott, 2015). Nevertheless, the older approach is still used extensively. It has been well understood for decades that the distributions of many response variables are not normal. Many are discrete (such as the counts of fruit or the percentages of diseased plants), while others are continuous but not symmetrical (possibly including the severity of disease symptoms or the time to an event, such as time to seed germination). The gamma and beta distributions are possible alternatives to the normal distribution for asymmetrical continuous distributions (Gbur et al., 2012; Ruíz et al., 2023). Counts with no (definable) upper bound (e.g., the number of insects on plants or the number of spores in a spore trap) may be described using the Poisson distribution, while counts with an upper bound [e.g., the number of plants with disease symptoms out of n plants observed (disease incidence) or the number of seeds germinating out of n seeds observed] may be described using a binomial distribution (Madden and Hughes, 1995; Madden et al., 2007; Gianinetti, 2020 ). In the binomial case, proportions can be defined as y/n, where y is the count. In the Poisson case, proportions do not exist as n is not defined. In practice, there is always a finite upper bound n for a count in biology; however, if n is very much larger than y, then it is reasonable to assume that there is no upper bound. An alternative for Poisson is the negative binomial (NB) distribution, while an alternative for the binomial is the beta-binomial distribution (Madden and Hughes, 1995). Both of these alternatives represent situations with higher variances than defined by the Poisson and the binomial. A typical property of non-normal distributions is that the variance is a function of the mean (Stroup, 2013). However, the standard assumption of a (normality-based) LM or LMM is that the (residual) variance is constant across all of the factor levels (i.e., the variance is independent of the mean). Linear or linear mixed modeling approaches can be (and routinely have been) used for non-normal data, as an approximation, by transforming the response variable so that the residual variance is roughly constant across the different means (e.g., for different treatments, cultivars, etc.). Examples include using the angular transformation (arcsine square root transformation) for proportions when the response variable has a binomial distribution and the square root transformation for unbounded counts when the response variable has a Poisson distribution (Piepho, 2003, 2009). Although this has been a popular and practical approach for decades, there are some disadvantages, as explained by Stroup (2015) and Gbur et al. (2012), among others. Essentially, this entails forcing the data to agree with a model (linear or LMM) rather than choosing a model that matches the stochastic process that generates the data (i.e., choosing a model that corresponds to the distribution of the data). Among other things, the estimated treatment “means” of the transformed values, or their back transformation to the original data scale, do not necessarily mean what users might think they mean. Additional issues are discussed further below. The analysis of non-normal data has been possible for decades using methods that are actually based on non-normal distributions, especially for data with binary, binomial, and Poisson distributions ( McCullagh and Nelder, 1989; Schabenberger and Pierce, 2002). Typical analyses include probit or logistic modeling of bioassay (binary dead or alive) data or count data. This was primarily for models with only fixed effects, with the models known as generalized linear models (GLMs), usually when the distribution belonged to the so-called exponential family of distributions. The addition of random effects to these GLMs to form GLMMs has posed some considerable statistical challenges, but research was in full force by the early 1990s (Breslow and Clayton, 1993; Wolfinger and O’Connell, 1993), although it would take another decade or more before relatively easy-to-use software became broadly available. The GLIMMIX procedure of SAS is especially important here, as is the lme4 package in R. GenStat also has methods for fitting GLMMs. One can consider LMs, LMMs, and GLMs as special cases of GLMMs. Readers should consult Littell et al. (2006); Gbur et al. (2012); Stroup (2013), and Ruíz et al. (2023) for more details on GLMMs, as well as on LMMs and GLMs. There are many reasons to use GLMMs in the agricultural and plant sciences (Gbur et al., 2012) because one can account for both random and fixed effects with several different realistic statistical distributions for the data. Nevertheless, there are also many issues that investigators need to be aware of when using GLMMs. These are described below after a more formal introduction of LMMs and GLMMs in the context of an example dataset. 3 Non-normal distributions 3.1 Example We consider a simple example to demonstrate several concepts and methods for the analysis of non-normal data. This is a subset of a much larger dataset analyzed in Paul et al. (2019) regarding, in part, the effect of fungicide treatments on Fusarium head blight in wheat (caused by the fungus Fusarium graminearum), mycotoxin contamination of the grain, and crop yield. This example is from one location–year (environment), a subset of the full dataset consisting of 29 environments (conducted by different researchers) with two factors evaluated at each environment, fungicide treatment, and cultivar resistance. We only consider the susceptible plant cultivar here, so that we are left with just one fixed-effects factor in a randomized complete block design (RCBD). If all the treatments were not in all of the blocks, we would have an incomplete block design. We are using this subset merely for demonstration purposes, not as a way of analyzing the full multi-environment dataset with two factors. The example experiment in this environment was laid out in four blocks (j = 1, ..., 4), with the six treatments (i = 1, ..., 6) randomized within each block. There were n = 100 wheat spikes randomly chosen and visually assessed for disease symptoms in each experimental unit (plot: a block–treatment combination identified by the ij index), and each spike was categorized as either diseased or healthy, giving a total of y[ij] diseased spikes per plot, with a proportion given by prop[ij] = y[ij]/n. This proportion is known as disease incidence (Madden et al., 2007), a measure of the fraction of individuals that are infected. With this design, there is a single response for each experimental unit (y[ij] or prop[ij]), even though this is the sum of several (n) binary observations. There is no requirement that n be constant as it could vary with plot (n[ij]). Here, we consider block to be a random effect and treatment to be a fixed effect. Note that, in Paul et al. (2019), more emphasis was placed on the analysis of the severity of disease on spikes (known as an “index” in the head blight literature) instead of the disease incidence. Severity is a continuous response variable representing the area of spikes with symptoms, expressed on a proportion or a percentage scale. The research questions are: Does treatment affect (mean) disease incidence? If so, which treatments are different from each other? Once we consider GLMMs for these data, we can ask better questions; for example, does treatment affect the probability of a wheat spike being infected (π) and which treatments differ from others in terms of π? To save space, the only difference we show is between the last and first treatments. 3.2 Distribution for disease incidence Since disease incidence is discrete with an upper bound of n, it is reasonable to consider the y in each plot to have a binomial distribution (Madden and Hughes, 1995), at least as a starting point. Generically, without reference to a particular treatment or block (i.e., no subscripts), we can write this as y ~ Bin(π,n), where π is a location parameter representing the probability of a trait or characteristic. For instance, π is the probability that a plant is infected or diseased; a moment estimate of π is the mean of the proportions for a given plot. The mean y for the binomial distribution is nπ and its variance is nπ(1 − π). Converting to proportions, the mean prop is π and its variance is π(1 − π)/n for the binomial. It is well known that the binomial can be well approximated using a normal distribution if n is large enough (Schabenberger and Pierce, 2002). To exemplify, following Stroup (2013, 2015), we simulated 100,000 observations from a binomial distribution with π = 0.1, with three different values of n (Figure 1). With n = 10, the binomial distribution is fairly skewed and poorly approximated by the normal (smooth curve); with n = 30, the binomial is much less skewed; and with n = 100, the binomial is very close to a normal distribution. With values of π closer to 0.5, the binomial is much closer to being symmetric (and is exactly symmetrical at π = 0.5), and approximation to normality is achieved with a smaller n. At very small or large π (i.e., very close to 0 or 1, respectively), a very large n may be needed to approximate Figure 1 Figure 1 Frequency distribution of 100,000 observations of the response variable y generated from a binomial distribution with n = 10, 30, 100 individuals per observation. Proportions were calculated as y/n for each observation. LMMs are not too sensitive to some departure from normality (Littell et al., 2006), and there are some LMM-fitting methods that do not require normality, e.g., MIVQUE0 (Rao, 1972), although calculations of confidence intervals and inference generally assume normality. However, even when it is reasonable to assume normality as an approximation for the binomial, the problem is that the variance of y (or of prop) will vary with the mean. The well-known angular (i.e., arcsine square root) transformation of prop does approximately stabilize variances, where prop* = sin^−1(√prop) ( Schabenberger and Pierce, 2002). Other variance-stabilizing transformations may be more successful when π is very close to 0 or 1 (Piepho, 2003). These may work well with LMMs (i.e., for assumed normal data), although it becomes non-trivial to interpret the exact meaning of the estimated angular means for each treatment relative to the underlying distribution of the counts (or corresponding proportions). This is especially true when there are random effects (see details and example below). There are also other challenges with the transformation-based LMM analysis. For instance, one transformation may be appropriate to stabilize variances, but a different transformation may be needed to obtain a linear (straight line) relation between y and a continuous explanatory variable (for a regression problem). Another transformation may be needed to obtain a symmetrical distribution. Thus, there are some big advantages to moving away from transformation-based analyses. The binomial is a member of the exponential family of distributions (McCullagh and Nelder, 1989; Gbur et al., 2012; Ruíz et al., 2023). There are several other distributions of the exponential family, or have statistical properties that are very similar to those of the members of the family. These include the Poisson, NB, gamma, and beta distributions (Stroup, 2013). The binomial, Poisson, and NB are for discrete data, whereas the gamma and beta are for continuous data. GLMs and GLMMs were developed to fit data from these distributions. The normal distribution, also known as the Gaussian, is a (symmetric) member of the exponential family, so LMM can be thought of as a special case of a GLMM. We focus on the normal and binomial distributions in this paper, which are applied to discrete data. Readers should consult Stroup (2013); Ruíz et al. (2023), and Gbur et al. (2012) for details on the other distributions. 4 Models and analysis 4.1 Linear mixed model Generically, we use y in the models as the response variable, with subscripts depending on the experimental and treatment design. For specific cases, we can use another symbol for the response. The classic LMM for a RCBD with one observation per experimental unit (ij combination) is: $\begin{array}{ll}{y}_{ij}=\text{θ}+{\tau }_{i}+{b}_{j}+{e}_{ij}& \left(1\right)\end{array}$ where θ is a constant (intercept), τ[i] is the effect of the i-th treatment (fixed), b[j] is the effect of the j-th block (random), and e[ij] is the residual (random), the latter representing variability in the response variable not accounted for by the other terms in the model. The residual essentially gives the unique effect of each experimental unit on the response variable (after accounting for the treatment and block main effects) and is equivalent to an interaction of the treatment and block effects (when there is one observation for each ij combination, as here). The random effects for this RCBD model are assumed to be normally distributed, with means of 0 and with variances ${\sigma }_{b}^{2}$ and ${\sigma }_{e}^{2}$: $\begin{array}{ll} {b}_{j}~N\left(0,{\sigma }_{b}^{2}\right),\text{\hspace{0.17em}}{e}_{ij}~N\left(0,{\sigma }_{e}^{2}\right)& \left(2\right)\end{array}$ Thus, Equations 1 and 2 are needed to jointly define the model for a RCBD with a random block effect. The (total) variance of y[ij] is ${\sigma }_{b}^{2}+{\sigma }_{e}^{2}$. In some circumstances, block could be considered as a fixed effect (Dixon, 2016), but we do not consider this here. Equation 1 (with the distributions in Equation 2) is defined as a mixed model because beyond the fixed effect(s), it has at least one random effect in addition to the residual (Stroup et al., 2018). It is linear (in terms of the parameters and response variable y) because the terms on the right-hand side of Equation 1 are actually shorthand expressions for a sum of parameters (factor effects) multiplied by binary indicator variables that identify treatments and blocks (Milliken and Johnson, 2009) and because the left-hand side does not involve any transformation of y. Note that the fixed effects are constants (to be estimated) and the random effects are (latent) random variables to be predicted (Stroup, 2013). Many statistical programs may refer to the random effect predictions as estimates and not predictions in the output. The non-residual random effects such as blocks do not need to be normally distributed (Lee and Nelder, 1996), but this is, by far, the most common assumption. Statistical research has shown that the results often are not very sensitive to this normality assumption for the random effects (McCulloch and Neuhaus, 2011; Schielzeth et al., 2020). It is a standard assumption of LMMs that the residuals have a normal distribution. The concept of a residual substantially changes when we transition to GLMMs. For a segue to GLMMs, Equation 1 can be rewritten in terms of expected values (means). That is, determining the expected value, E(·), of the left- and right-hand sides of Equation 1 conditional on the block random effect (more generally, conditional on all random effects other than the residual) leads to: $\begin{array}{ll}{\mu }_{ij}=\text{θ}+{\tau }_{i}+{b}_{j}& \left(3\right)\end{array}$ where μ[ij] is the so-called conditional expected value or the mean of the response for treatment i and block j. It is “conditional” because its value is for the specific level(s) of the random effect (the j-th block for this simple RCBD). From Equation 3 and the normality assumptions, the distribution of the response variable y[ij] conditional on the random effect, y[ij]|b[j], is defined $\begin{array}{ll}{y}_{ij}|{b}_{j}~N\left({\mu }_{ij},{\sigma }_{e}^{2}\right)& \left(4\right)\end{array}$ Note that the mean for this conditional distribution of y[ij] is μ[ij] and that the variance of this conditional distribution for a given block (or a given level of the random effects) is ${\sigma }_ {e}^{2}$ (which was given as the residual variance in Equation 2). In other words, after accounting for the fixed and random effects, the only other variability of y is captured by the variance of the conditional distribution (i.e., the unexplained variability). Readers should note that the μ[ij] in Equations 3 and 4 may be written as μ[ij]|b[j] in order to make the conditioning explicit (we avoid this additional notation here). Putting it all together, the alternative to Equations 1 and 2 is: $\begin{array}{ll}\begin{array}{l}{\mu }_{ij}=\text{θ}+{\tau }_{i}+{b}_{j}\\ {y}_{ij}|{b}_{j} ~N\left({\mu }_{ij},{\sigma }_{e}^{2}\right)\\ {b}_{j} ~N\left(0,{\sigma }_{b}^{2}\right)\end{array}& \ In this manner of writing the LMM, the residual term “disappears”; ${\sigma }_{e}^{2}$ is now seen as the variance of the conditional normal distribution for y (after accounting for other effects). Equation 5 can actually be generalized further in a way that will help to understand GLMMs for non-normal response variables. The right-hand side of Equation 3 for the conditional mean is known as the linear predictor (LP), which is typically written with the η symbol (η[ij] for the RCBD here). The LP specifies all the fixed and random effects in a controlled experiment or with observational data that can affect either the conditional mean of a response variable (or a function of the conditional mean). For a LMM (i.e., normal conditional distribution), the mean simply equals the LP; that is, we link the LP to the mean with ${\mu }_{ij}={\eta }_{ij}$. We can then rewrite the set of equations for a RCBD as: $\begin{array}{ll}\begin{array}{l} {\eta }_{ij}=\text{θ}+{\tau }_{i}+{b}_{j}\\ {\mu }_{ij}= {\eta }_{ij}\\ {y}_{ij}|{b}_{j} ~N\left({\mu }_{ij},{\sigma }_{e}^{2}\right)\\ {b}_{j} ~N\left(0,{\sigma }_ {b}^{2}\right)\end{array}& \left(6\right)\end{array}$ The LMM is rarely seen written this way, and the inclusion of the explicit LP component is excessive for this LM with conditional normal distribution; however, it does nicely allow for a transition to the non-normal conditional distributions of GLMMs. This formulation is also very useful for Bayesian model fitting because it defines the two statistical distributions in the model: the distribution of the block random effect and the conditional distribution of the response variable. Equations 1 and 2, or Equation 6 (with the four sub-equations), can be expanded for any number of fixed and random effects, as well as the interactions of random and fixed effects (Littell et al., 2006; Stroup, 2013). Both formulations are conditional models, with the same results when fitted to data; it is just a matter of preference in how one writes the LMM. From Equations 5 or 6, the expected mean for the i-th level of the fixed effect (the i-th treatment) can be determined. The random effects are integrated out to obtain μ[i] (instead of μ[ij]). Because the expected value of the normally distributed b[j] is 0, the result for a LMM is simply: $\begin{array}{ll}{\mu }_{i}=\text{θ}+{\tau }_{i}& \left(7\right)\end{array}$ This is generally of most interest to the investigator. The estimates of μ[i] and the differences of μ[i] between treatments (μ[i] − μ[i][′], where i and i′ are any two treatments), as well as their SEs (a function of the variances in Equation 6), are derived from statistical theory and are automatically calculated with LMM software (with the right statements or options). Details are in mixed-model textbooks (e.g., Littell et al., 2006; Galecki and Burzykowski, 2013; Stroup et al., 2018). There is yet another way to write a LMM. The two random-effects terms (including the residual) in Equation 1 can be taken and combined into one term as: h[ij] = b[j] + e[ij]. The LMM can then be written as: $\begin{array}{ll}{y}_{ij}=\text{θ}+{\tau }_{i}+{h}_{ij}& \left(8\right)\end{array}$ This LMM is known as a marginal model for a RCBD experiment (Stroup, 2013). Both block and residual effects are still present, just expressed differently. The variability component(s) now has to be written with more complexity to express both random sources. The approach is to define a vector of random effects for each level of the random effect (the j-th block here), designated h[j]. In order to save space, we assume here that there are just three treatments, so we can write: $\begin{array}{ll}{\mathbf{\text{h}}}_{\mathbf{\text{j}}}=\left(\begin{array}{c}{h}_{1j}\\ {h}_{2j}\\ {h}_{3j}\end{array}\right)=\left(\begin{array}{c}{b}_{j}+{e}_{1j}\\ {b}_{j}+{e}_{2j}\\ {b}_{j}+ {e}_{3j}\end{array}\right)& \left(9\right)\end{array}$ where each element of the vector corresponds to a different treatment. This leads to the derivation of the variance–covariance matrix (P) for y[ij]: $\begin{array}{ll}P=\left(\begin{array}{ccc}{\sigma }_{\text{b}}^{2}+{\sigma }_{\text{e}}^{2}& {\sigma }_{\text{b}}^{2}& {\sigma }_{\text{b}}^{2}\\ {\sigma }_{\text{b}}^{2}& {\sigma }_{\text{b}}^{2}+ {\sigma }_{\text{e}}^{2}& {\sigma }_{\text{b}}^{2}\\ {\sigma }_{\text{b}}^{2}& {\sigma }_{\text{b}}^{2}& {\sigma }_{\text{b}}^{2}+{\sigma }_{\text{e}}^{2}\end{array}\right)& \left(10\right)\end The diagonal elements of the matrix are the (total) variances of y[ij] (${\sigma }_{b}^{2}+{\sigma }_{e}^{2}$; the same for each treatment), while the off-diagonal elements (${\sigma }_{b}^{2}$) are the covariances of the response variables. This formulation helps to make clear that, with a random block, observations that are from the same block are correlated. The correlations within a block can be determined from the covariances and variances (Littell et al., 2006). This structure is known as compound symmetry (CS). The conditional model (Equations 1 and 2, or Equations 5 or 6) and the marginal model (Equations 8 and 10) are equivalent for LMMs. Fitting either version will give the same result when fitted to data if ${\sigma }_{b}^{2}$ is 0 or positive. Although the conditional model in Equations 5 or 6 does not explicitly show a covariance, the covariances between treatments are still there, derivable from the terms of the conditional model. For the RCBD, the correlation of y[ij] within the same block (say, j = 1), known as the intra-class correlation, is: $\begin{array}{ll}\rho ={\sigma }_{b}^{2}/\left({\sigma }_{b}^{2}+{\sigma }_{e}^{2}\right)& \left(11\right)\end{array}$ The reader can find more details in Stroup (2013) and Gbur et al. (2012). The most common method to fit LMMs is through the use of restricted maximum likelihood (REML), also known as residual maximum likelihood. This can be considered the gold standard. REML produces unbiased parameter estimates (such as with a RCBD) or less biased estimates than those produced with maximum likelihood (ML). Bayesian methods can also be used (Wolfinger and Kass, 2000; Stroup, 2021 ), as well as some alternative frequentist approaches (Littell et al., 2006). It is important to note that, whether the data analyst uses the conditional (Equation 6) or the marginal (Equations 8-10) form of the LMM, the actual (observed) data correspond to the marginal distribution (Gbur et al., 2012). In essence, the conditional and marginal models are just two different ways of generating data from a marginal distribution (Stroup, 2013). The random effects in a conditional LMM are random latent variables that are not directly measured or observed; rather, they are predicted through the model-fitting procedure. What one observes (measures) and then analyzes are data arising from the marginal distribution, no matter how that marginal distribution is derived or generated. 4.2 Example data analyzed with LMMs Using the GLIMMIX procedure in SAS with REML estimation, the conditional LMM in Equation 6 (equivalent to Equations 1 and 2) was fitted to the example data with the proportion of diseased wheat spikes as the response variable (prop[ij] for y[ij]). This is not normally done as the variance is not independent of the mean for binomial data (or discrete data, in general). Estimates of the treatment means (e.g., $\stackrel{^}{\theta }+{\stackrel{^}{\tau }}_{1}={\stackrel{^}{\mu }}_{1}$) and the SEs are given in Table 1, together with one pairwise difference of means. The exact same results are obtained if the marginal model (Equations 7–10) is fitted for this LMM (results not shown). The SEs are functions of the block and residual variances (details in Gbur et al., 2012; Stroup, 2013). Interestingly, despite the problems in directly analyzing proportions, the estimated mean proportions from the fitted LMM are unbiased estimates of the true proportions for each treatment across all the blocks (i.e., across all the levels of the random effects; the marginal means) (Stroup, 2013, 2015). However, the estimated SEs across all treatments (0.0590) are meaningless because the LMM assumes a constant variance (${\sigma }_{e}^{2}$). Since the estimates of variability are wrong, the estimates of confidence intervals and the tests of significance are thus wrong. Table 1 Table 1 Estimated means and standard errors (SEs) when fitting a linear mixed model (LMM; Equation 6) to the proportion of diseased wheat spikes^a with an incorrect assumption that the residual variance is constant (i.e., independent of the mean proportion) and when fitting the LMM to the angular transformation of the proportions where it is reasonable to assume the transformed values have constant residual variance, together with the back-transformation of the estimated angular means to obtain proportions (where the SE of the back-transformation is based on the delta method). The LMM of Equation 6 can be generalized to allow for a separate residual for each level of the fixed effect (${\sigma }_{ei}^{2}$). This is easily done using REML algorithms such as in the MIXED procedure in SAS. Fitting such a model can be more challenging, especially with the small sample sizes typically found in agricultural field studies. With the number of blocks in this example being four, attempting to estimate each treatment-specific residual variance is thus not a good strategy. The results of fitting the LMM to the transformed angular data are also shown in Table 1. The estimated SEs were constant across the treatments (0.0651), which is expected since one assumes that the angular-transformed data have variances independent of the mean. Comparisons of the treatments and tests of significance can be done on these mean transformed values, and back-transformation of the means can be performed to obtain a type of “mean proportion” for each treatment, as shown in the table. However, it is not intuitive what these back-transformed means represent exactly. They are not estimates of the means of the distributions of the actual proportions, but are related to those means. In symbols and suppressing subscripts here, using g(y) as a transformation (function) of y [such as a log and angular (arcsine square root), among others], the expected value (i.e., mean) of g(y), E( g(y)), is not equal to the transformation of the expected value (mean) of y, g(E(y)) (Schabenberger and Pierce, 2002). These two types of means are related, but are not the same. The situation becomes even more complicated with multiple random effects. For the example (Table 1), the back-transformed values give a type of “central” (but not truly central) value of the marginal distribution of the data. Nevertheless, analytical approaches based on transformations still have value, as long as the investigators understand the limitations. 5 Generalized linear mixed models 5.1 Conditional non-normal distributions With GLMMs, investigators have the opportunity to better match the model used for data analysis with a plausible model for stochastic generation of the data (Stroup, 2013, 2015). Consider the disease incidence data in the example, where it is reasonable to assume a binomial distribution for y in a given plot (experimental unit), at least as a starting point. Suppressing subscripts temporarily, we can write generically, y|plot ~ Bin(π,n), for the conditional distribution of the number of diseased wheat spikes (this could also be for the number of germinated seeds in a pot in a greenhouse, etc.). The probability of disease (or the probability of some trait), π, could be influenced by the pathogen inoculum density, the favorability of the environment for infection, and so on. It is quite plausible that π represents a parameter that has direct mechanistic interpretation. We continue with our hypothetical example from Figure 1 and consider π = 0.1 (following closely, once again, the method in Stroup, 2013). Consider the impact of random effects, such as blocks. If there is only one treatment (e.g., control), then the plot and the block are synonymous. With random effects, π is randomly perturbed by the block, so that π is higher than 0.1 in some blocks (plots) and is lower than 0.1 in other blocks. The response y in each block (for any single treatment) would depend on the realized value of π for that block. On average, the perturbation is 0, at least on one possible scale. Ultimately, we want to specify the marginal distribution of y [across the population of plots (blocks)] based on the data-generating process in each plot (block) and the distribution of perturbations across the plots. A linear additive scale for the perturbation of the random block effects (i.e., adding b[j] directly to π) would not work in a realistic model, in general, because a random permutation could give a probability nonsensically above 1 or below 0 in a given plot when π is relatively close to 1 or 0. A linear perturbation could work on a function of π, if the right function was used, so that π remains bounded by 0 and 1. This is the approach taken in the usual GLMM. Being more explicit and referring back to the conditional distribution form of a LMM (Equation 6), we can write the LP for a RCBD as $\text{ }{\eta }_{ij}=\theta +{\tau }_{i}+{b}_{j}$. This is exactly the same for a normal or a non-normal conditional distribution. Likewise, the random block effect can be assumed to have a normal distribution, as with LMMs. The question is how to link the η [ij] LP to the location parameter. This is done through the link function, g(π[ij]), or, more generally, as g(μ[ij]), where a link function is a transformation of a parameter (not a transformation of the observations). Thus, we can write g(π[ij]) = η[ij]. The location parameter is obtained by the inverse link function, written as π[ij] = g^−1(η[ij]) = h(η[ij]). For conditional binomial data, the logit is the most common link function; that is, g(π[ij]) = logit(π[ij]) = ln(π[ij]/(1 − π[ij])). There are other possibilities for the link function based on specific applications or theory ( Collett, 2003). The term “identity link” or “identity link function” is used when there is no transformation of μ or π. For a RCBD (one fixed-effects factor and random block effect factor), we can now write the conditional model form of the GLMM generically as: $\begin{array}{ll}\begin{array}{l} {\eta }_{ij}=\text{θ}+{\tau }_{i}+{b}_{j}\\ g\left({\mu }_{ij}\right)= {\eta }_{ij}\\ {y}_{ij}|{b}_{j} ~\text{Dist}\left({\mu }_{ij},n\right)\\ {b}_{j} ~N\left(0,{\ sigma }_{b}^{2}\right)\end{array}& \left(12\right)\end{array}$ where “Dist” is generic for some conditional distribution. This equation is almost the same as that for the LMM (see Equation 6). For conditional binomial data (with π = μ), in particular, we write: $\begin{array}{ll}\begin{array}{l}\text{ }{\eta }_{ij}=\text{θ}+{\tau }_{i}+{b}_{j}\\ \text{logit}\left({\pi }_{ij}\right)=\text{ }{\eta }_{ij}\\ {y}_{ij}|{b}_{j} ~Bin\left({\pi }_{ij},n\right)\\ {b} _{j} ~N\left(0,{\sigma }_{b}^{2}\right)\end{array}& \left(13\right)\end{array}$ Recall that the mean of the binomial is nπ for y; the mean of y/n is π. Thus, the location parameter π with the binomial plays the role of μ in the general distribution (Equation 12). Unlike with a LMM, it is important to note that there is no variance to estimate for the conditional binomial distribution (i.e., no unknown residual variance as the variance of the binomial is fully defined based on n and π). The variance of the conditional binomial is nπ(1 − π). See GLMM expansions and extensions below (Section 6) for the implications of this. One can reduce Equation 13 to three components by directly writing $\text{ logit}\left({\pi }_{ij}\right)=\text{θ}+{\tau }_{i}+{b}_{j}$ as the first sub-equation. It is important to note that, although the g(·) link function is a transformation, a GLMM is not a model for a function or transformation of data (i.e., not a model for the transformed response variable). Put another way, a GLMM is not a model for the mean of transformed observations; rather, it is a model for a function (transformation) of the mean (μ, π, etc.) of the actual (non-transformed) observations. This subtle difference is extremely important. The LP of the GLMM in Equation 13 estimates a function of the probability of disease for a particular treatment and block. The inverse link to obtain π for a logit link function is given by: $\begin{array}{ll}{\pi }_{ij}=h\left({\eta }_{\text{i}j}\right)=1/\left(1+\mathrm{exp}\left(-\text{logit}\left({\eta }_{ij}\right)\right)\right)& \left(14\right)\end{array}$ 5.2 Marginal distribution Integrating over the random effects of the conditional LMM (Equation 6) results in the marginal distribution for y (Littell et al., 2006; Stroup, 2013; Brown and Prescott, 2015). With normality for the random effects and the conditional distribution in a LMM, the marginal distribution is determined analytically. The expected values of the fixed effects of this marginal distribution (e.g., $\ text{ }{\mu }_{ij}=\text{θ}+{\text{τ}}_{\text{i}}$) are the same as those obtained simply by inserting 0 (the expected value of the normally distributed random effects; see Equation 2) for the random effects in the conditional LMM (Equation 3). The situation is different with GLMMs. The marginal distribution cannot be determined analytically, requiring numerical integration; for the binomial conditional distribution, as an example, $\text{ logit}\left({\pi }_{i}\right)=\text{θ}+{\tau }_{i}$ is not the logit of the mean proportion of diseased individuals over all the blocks (i.e., over all the levels of the random effects). This is because the inverse link is a nonlinear equation (Equation 14) and the conditional distribution is not normal. This is an extremely important concept because, as discussed above, one observes the marginal distribution (more technically, a sample from the marginal distribution), not the conditional distribution. Numerical integration over the random effects can give an estimate of the marginal distribution (required for model fitting), including estimates of the marginal mean (the marginal mean proportion of diseased plants here over all the blocks) and its variability. This is explained by Stroup (2013) on pages 102–107 for a similar situation. Sometimes, this marginal distribution, which is made up of the mixture of the binomial and normal distributions, is called the logistic–normal–binomial (LNB) distribution when used to characterize the spatial variability in fields (Hughes et al., 1998). The marginal distribution from Equation 13 can also be obtained by simulation for known parameters for the conditional distribution and the random-effects distributions. Following Stroup (2013) closely, we estimated the marginal distribution of the diseased proportion (y/n) across levels of the random effects by simulating 100,000 observations from Equation 13, with π = 0.1 and n = 30 for the conditional binomial distribution, and three levels of the block effect standard deviation [${\sigma }_{b}$ = 0.5, 1.0, and 1.5 (${\sigma }_{b}^{2}$ = 0.25, 1.0, and 2.25)]. With simple random sampling (i.e., no random effect distributions), the binomial distribution is not very skewed with n = 30 (Figure 1). However, in a mixed-model setting with symmetric (normal) random effects, the marginal distribution of the proportion can be very skewed depending on the value of the block variance. As seen in the left-hand column of Figure 2, the marginal distribution becomes more and more skewed with increasing block variance. A general rule of thumb is that marginal distributions are not symmetrical for GLMMs, with the skewness dependent on the skewness of the conditional distribution (e.g., binomial) and the magnitude of the variances of all the random effects in a LP. Figure 2 Figure 2 Marginal distributions determined by generating 100,000 observations of the response variable y with a conditional binomial distribution (with π = 0.1 and n = 30 individuals per observation), plus a random effect for each observation (which had, for most graphs, a normal distribution with variance of ${\sigma }_{b}^{2}$or standard deviation of ${\sigma }_{b}$). Inverse link is used to obtain y from logit(π). Proportion is calculated as y/n. Left-hand column: (A–C) Data generated with Equation 13 (without τ because there is only one treatment, with b[j] now as a plot effect), with standard deviations of b[j] effect of 0.5 (A), 1.0 (B), and 1.5 (C). Right-hand column: (D) Data generated as in (B) (standard deviation of 1), but the proportions transformed with angular transformation. (E) Data generated as in (B) (standard deviation of 1), but the proportions transformed with Piepho (2003) function (with δ = 5). (F) Unlike other graphs, y is generated with beta-binomial distribution with overdispersion parameter of $\varrho$ = 0.1 and no b[j] effect, approximately giving the same variance of y/n as in (B) (with a standard deviation of 1). In the simulation example in Figure 2, the mean proportions of the marginal distributions are estimated as 0.109, 0.134, 0.169 for block variances of 0.25, 1.0, and 2.25, respectively, all with π = 0.1 (displayed in Figures 2A–C, respectively). That is, the mean proportions increase with increasing random effect variances even though the underlying location parameter for the conditional distribution (i.e., for the generation of observations in each experimental unit) is fixed at π. Broadly, the mean proportion will not correspond to the inverse link of θ + τ[i] for treatment i (or just θ when there is just one treatment). When π is less than 0.5, the marginal mean proportion will be greater than π, with the degree of difference being dependent on the random effect variability. When π is greater than 0.5, the marginal mean proportion will be less than π. The variances of the marginal distributions also increase in the same manner. Using the angular transformation of the observations in a LMM does not lead to a symmetric marginal distribution, as seen in Figure 2D for an example block variance of 1.0. The new variance-stabilizing transformation of Piepho (2003), with a stabilizing parameter of δ = 5 also does not lead to a symmetric distribution when the block variance is 1.0 in this example (Figure 2E). We estimated δ using the ML methodology in Piepho (2003) (results not shown). Other tested values of δ did not result in a symmetric distribution. As discussed by Piepho (2003), the new variance-stabilizing transformation to be used in LMMs is very effective in stabilizing variances (among treatments), but is less effective in producing a symmetrical distribution of the transformed data (especially when the original distribution of y is very There are other ways to derive a marginal distribution than the typical method here of mixing a normal random effect distribution with conditional binomial on the link scale. For example, it is well known that the beta-binomial provides an excellent representation of the distribution of disease incidence when there is overdispersion (i.e., higher variability than for the binomial alone, sometimes called extra-binomial variation) (Madden and Hughes, 1995; Turechek and Madden, 1999; Madden et al., 2007). In general, the beta-binomial is utilized to characterize the spatial variability (clustering) of incidence within fields, where each sampling location in a field is a separate observation of disease incidence (based on n observed units for each sample) (Turechek and Madden, 1999 ), but the concept is easily extended to variability among blocks in planned experiments. In the context of a RCBD, the beta-binomial is derived by, once again, assuming that π is randomly perturbed in each plot (block, etc.). However, the random block effect is not additive on the link scale of the binomial [e.g., not adding b[j] to the equation for logit(π[j])] as it is for the derived GLMM of Equation 13. Instead, the random block effect is assumed to have a beta distribution and is multiplicative with π (i.e., working directly on the scale of π) (Madden and Hughes, 1995). The plot effect is characterized by a $\varrho$ correlation parameter (instead of ${\sigma }_{b}^{2}$). An example of simulated observations from a beta-binomial distribution is shown in Figure 2F, where $\varrho$ was chosen to roughly have the same effect as ${\sigma }_{b}^{2}$ = 1 with the normal random effect (Figure 2B). As with the other scenarios in Figure 2, the marginal distribution is very skewed There are many uses of the beta-binomial and the related binary power law (Hughes and Madden, 1995; Madden et al., 2018), but the beta-binomial is not commonly used for routine data analysis with GLMMs (e.g., determining treatment effects). Since the random effect is considered non-normal, the beta-binomial is a special case of what is often called doubly (or double) generalized linear mixed models (DGLMMs) (Lee et al., 2006). We consider this in some more detail below. The key question in a GLMM is the measure of the “location” (and associated SE) one wants to estimate. In agreement with Stroup (2013, 2015) and Stroup et al. (2018), we argue that for conditional binomial data, one wants to estimate logit(π[i]) = θ + τ[i] for treatment i and the conditional probability π[i] using the inverse link function (location parameter of the conditional binomial) ( Equation 14). That is, one is interested in the parameter from the conditional distribution, the probability of disease when the random effects are exactly at their means (0). This is considered a conditional mean. This location parameter is not affected by the distribution of the random effects [although the random effect variances do affect parameter estimation and the precision (SE) of parameter estimates]. It turns out that the π[i] from the conditional binomial (for treatment i) is the median of the marginal distribution of proportions for this treatment (Stroup, 2013). Most importantly, when the GLMM of Equation 13 is fitted to RCBD data, estimates of logit(π[i]) for each treatment (θ + τ[i]), and therefore π[i] through the inverse link (Equation 14), can be directly 5.3 Example The estimates of the logit of the probability of disease and the corresponding conditional probabilities (inverse links), together with their estimated SEs, from the fit of Equation 13 are given in Table 2 (listed as model A). This is based on restricted pseudo-likelihood (RSPL) model fitting, discussed below. As expected, the SEs are not equal for the different treatments, being a function of the variance of the conditional binomial and the block variance. The block variance estimate was small in this example (${\stackrel{^}{\sigma }}_{b}^{2}$ = 0.073), and the proportions are not very close to 0 or 1, meaning that the ${\stackrel{^}{\pi }}_{i}$ of the conditional binomial (based on the inverse link) is close to the mean proportion of the marginal distribution obtained with the fit of the LMM to the proportion data (Table 1). However, the SEs vary greatly between the LMM and GLMM approaches [e.g., 0.042 for treatment 1 with the GLMM (model A in Table 2) versus 0.059 for the same treatment with the LMM (Table 1)], reflecting that the GLMM properly accounted for the variance of a binomial distribution and estimated a conditional distribution parameter. Table 2 Table 2 Estimated means and standard errors (SEs, in parentheses) both on the linear predictor scale [logit(mean)] and the data scale (inverse link; with SEs based on the delta method) when fitting three generalized linear mixed models (GLMMs) using restricted pseudo-likelihood (models A and B) or restricted pseudo-likelihood with a quasi-binomial conditional likelihood (model C) to the number of diseased wheat spikes^a: naive Equation 13 (model A); Equation 16, which includes an additive random effect for the individual experiment unit (plot; model B); and Equation 17, which includes a multiplicative overdispersion parameter instead of an additive experimental unit effect (model C). 6 GLMM expansions and alternatives 6.1 Link functions Although it is very common to use the logit link with a conditional binomial distribution in a GLMM, there are other choices, such as the probit (inverse standard normal function) and complementary log–log (CLL) [ln(−ln(1 − π))]. There may be theoretical reasons to choose a different link for some applications (Collett, 2003), or selection may be based on the goodness of fit of the model (Malik et al., 2020). For instance, Kriss et al. (2012) used the CLL link in a hierarchical GLMM for surveys of clustered plant disease incidence data, where the choice of link was based on previous theoretical and empirical work on the relationship between disease incidence and severity (a continuous variable) and the relationship between incidence at two (or more) scales in a spatial hierarchy (Turechek and Madden, 2003; Hughes et al., 2004; Paul et al., 2005). Many other possible link functions have been proposed, and Malik et al. (2020) described several of them, with a proposed approach for deciding on an appropriate link to use. The reason for the logit being the default link has to do with the form of the binomial distribution. One method of writing the log of the binomial distribution (derived after a little algebraic manipulation of the function typically given in introductory statistics courses) is: $\begin{array}{ll}\text{ln}\left[f\left(y\right)\right]=y \text{ln}\left(\frac{\pi }{1-\pi }\right)+n\mathrm{ln}\left(1-\pi \right)+ln\left(\begin{array}{c}n\\ y\end{array}\right)& \left(15\right)\ Equation 15 is in the form for a member of the exponential family of distributions (Collett, 2003; Gbur et al., 2012). A key property of this family of distributions is that there is a term that involves the product of the response variable (y) and a so-called canonical parameter, which is either the mean (expected value, μ) or the location parameter, or a function of the mean for the distribution. Here, the canonical parameter is the logit of π, ln(π/(1 − π)), the typical link function for the binomial. For the Poisson, for instance, the canonical parameter is ln(μ), which is the typical default link function for count data. The work by Gbur et al. (2012) and many of the major reference books for GLMMs give the ln[f(y)] forms for other distributions or density functions in the exponential family. All the relevant properties of these distributions (such as the variance, among others) follow directly from the form of the ln[f(y)]. When the default canonical link is not selected, the main software programs, such as GLIMMIX in SAS, handle the necessary extra calculations in the background to fit models and calculate statistics. 6.2 Realistic modifications of the GLMM The GLMM with a binomial conditional distribution (Equation 13) is the obvious natural extension of the LMM (Equation 6) for a RCBD; however, there is a very good chance that the use of the simple GLMM with agricultural (or other) experiments will give incorrect results! In particular, the SEs of the estimated location parameters (μ and π) for treatments (on the link scale of the model or for the inverse link scale of the data) may be too small. Likewise, tests of significance are likely to be wrong or misleading, where the test statistics are too large and the p-values for significance too small (Stroup, 2015). This is an important warning, as elaborated below. To understand this, note that the conditional normal distribution in Equation 6 has a variance (${\sigma }_{e}^{2}$) that is estimated, but there is no variance to estimate with the conditional binomial. The variance of y for a conditional binomial distribution is fully defined by π and n; that is, var(y[ij]|b[j]) = nπ [ij](1 − π[ij]) for the RCBD with random block effect. However, as discussed elsewhere, count data are very commonly overdispersed (Madden et al., 2007, 2018); that is, y has a higher variance at the experimental unit (plot) level than nπ[ij](1 − π[ij]). Similarly, for counts without an upper bound, y has a higher variance at the plot level than the theoretical variance for a Poisson distribution (μ) (Madden and Hughes, 1995; Madden et al., 2018). As an aside, the concept of overdispersion as used here does not exist for the (conditional) normal distribution because the residual variance is independent of the mean; ${\sigma }_{e}^{2}$ can take on any value to account for the variability not accounted for by the other terms in the model (in other words, the residual variance is estimated from the data). For plant diseases, there are numerous causes for overdispersion at the experimental unit (plot in this case) level, especially the clustering of diseased individuals in plots or fields (Hughes and Madden, 1995; Madden et al., 2007). More broadly, any non-uniformity within the experimental units or nonrandom sampling of the units, or misspecification of the LP (right-hand side of the equation for η), or misspecification of the conditional distribution in the model can result in overdispersion in a given data analysis. Put another way, overdispersion can occur “when the model fails to adequately account for all the sources of variability” (Stroup et al., 2018, p. 391). There are multiple ways to deal with overdispersion when analyzing discrete data. The primary approaches for our purposes in this article are: 1) to expand the LP to include an additional random-effects term (or terms); 2) to use a quasi-likelihood for the conditional distribution; or 3) to use a different conditional distribution that allows for higher variability than for the binomial (or Poisson for unbounded counts) (Stroup, 2013). There are important implications for all of these model expansions in terms of model fitting (estimation). We defer the discussion on estimation until later. 6.3 Adding terms to the linear predictor The first of these approaches is the expansion of the GLMM. For the RCBD conditional model of Equation 13, we add a term for the interaction of block and treatment, (bτ)[ij], which we can write as v [ij], to the right-hand side of the LP: $\text{ }{\eta }_{ij}=\text{θ}+{\tau }_{i}+{b}_{j}+{v}_{ij}$ The distribution of v[ij] is assumed to be normal, ${v}_{ij}~N\left(0,{\sigma }_{v}^{2}\right)$. Now, the conditional distribution of y given the random effects is written as: ${y}_{ij}|{b}_{j},{v}_{ij} ~Bin\left({\pi }_{ij},n\right)$ Conceptually, we assume that π is randomly perturbed not only by the block but also by the individual plot. Recall that the residual term in Equation 1 for a normality-based LMM is equivalent to the interaction of block and treatment, where the residual term accounts for the variability not accounted for by the other terms in the model. Therefore, the expansion of the LP for the conditionally binomial data is in the same spirit, taking into account the experimental design [for experiments without blocks, the b[j] term would not be present, but the v[ij] term would still be used here for the random effect of the experimental unit (plot)]. For the RCBD with conditional binomial data, Equation 13 is now written as: $\begin{array}{ll}\begin{array}{l}\text{ }{\eta }_{ij}=\text{θ}+{\tau }_{i}+{b}_{j}+{v}_{ij}\\ \text{logit}\left({\pi }_{ij}\right)=\text{ }{\eta }_{ij}\\ {y}_{ij}|{b}_{j},{v}_{ij} ~Bin\left({\pi }_ {ij},n\right)\\ {b}_{j} ~N\left(0,{\sigma }_{b}^{2}\right)\\ {v}_{ij} ~N\left(0,{\sigma }_{v}^{2}\right)\end{array}& \left(16\right)\end{array}$ The same general approach is used for any conditional GLMM that is based on a conditional binomial or a conditional Poisson distribution (with any number of fixed and random effects), i.e., to add a residual-like random effect to the LP equation of the model. This may be a three-way or a higher-order interaction, generally a crossing of all the terms in the LP. Equation 16 is considered a conditional model, and the π[ij] represents the probabilities of disease for the individual plots (combination of treatment and block), the same as with Equation 13. Interest still remains on the estimation of logit(π[i]) (= $\text{θ}+{\tau }_{i}$) and, hence, π[i]. 6.3.1 Example Fitting Equation 16 using RSPL (see below for a discussion on model fitting), which we label as model B, resulted in estimates of ${\stackrel{^}{\sigma }}_{b}^{2}$ = 0.031 for block (compared with 0.073 for model A) and ${\stackrel{^}{\sigma }}_{v}^{2}$ = 0.248 for the block × treatment interaction [plot random effect (“residual”)] (Table 2). Although the block variance is smaller than that of the naive GLMM above (no v[ij] term), there is now a nonzero plot variance that was not previously considered. The estimates of the LP (logit) and the inverse link for each treatment are shown in Table 2. Here, the estimated SEs are larger than that for the fit of the simpler Equation 13 (e.g., 0.283 for treatment 1 logit for Equation 16 versus 0.168 for Equation 13). This, together with the nonzero estimated variance for v[ij], is an indicator of the overdispersion that needed to be accounted for. The F statistic for treatment significance is now 5.24, smaller than that for the simpler (and inadequate) Equation 13 that did not account for the extra variability (F = 27.9). The means (on the link or data scale) are similar to the results found using the simpler model. However, these means will, in general, not be the same. 6.4 Quasi-likelihood When there is greater variability than is possible with a binomial distribution (i.e., with overdispersion or extra-binomial variation), instead of adding a random effect term as the LP, the conditional variance can be defined as a constant times the binomial variance. For a RCBD, this is var(y[ij]|b[j]) = ϕnπ[ij](1 − π[ij]), where ϕ is an overdispersion scale parameter. It plays the role of the index of dispersion (D) in spatial pattern analysis (Madden et al., 2018). By allowing ϕ to take on any positive value, the conditional variance can always be expressed as a multiple of the binomial. For unbounded counts, one would write var(y[ij]|b[j]) = ϕμ[ij], where μ[ij] is the variance (and the mean) of the Poisson distribution. With ϕ not equal to 1, the conditional distribution is no longer binomial; in fact, there is no longer an actual true conditional statistical distribution that could stochastically generate the data. Thus, one has a so-called quasi-likelihood, and a method suitable for quasi-likelihood is used to fit the model (Gbur et al., 2012; Stroup, 2013). See below for a discussion on estimation. The quasi-likelihood-based model for the RCBD is given as a slight modification of Equation 13: $\begin{array}{ll}\begin{array}{l} {\eta }_{ij}=\text{θ}+{\tau }_{i}+{b}_{j}\\ \text{logit}\left({\pi }_{ij}\right)=\text{ }{\eta }_{ij}\\ {y}_{ij}|{b}_{j} ~ quasi.Bin\left({\pi }_{ij},n;\varphi \ right)\\ {b}_{j} ~N\left(0,{\sigma }_{b}^{2}\right)\end{array}& \left(17\right)\end{array}$ where $quasi.Bin\left({\pi }_{ij},n;\varphi \right)$ is an arbitrary expression that represents a conditional quasi-likelihood (which is not for the actual distribution of the data). Note that the LP (η[ij]) is the same as the naive GLMM in Equation 13 (i.e., no addition of a v[ij] term). It would be incorrect to add the v[ij] term (block × treatment interaction) in a model where the ϕ is also added as this would be overparameterization; that is, two terms (${\sigma }_{v}^{2}$ and $\varphi$) would be “competing” to explain the overdispersion. 6.4.1 Example Table 2 displays the estimated means and SEs (model-scale logits and data-scale inverse link means) for the fit of Equation 17 based on the conditional quasi-binomial likelihood (model C). This was done with RSPL (see below). The estimated means are extremely close to those obtained for the fit of naive Equation 13 (model A), but the SEs are much larger, reflecting an estimated scale parameter of $\stackrel{^}{\varphi }$ = 5.77 (i.e., the conditional variance is nearly six times larger than that obtainable with the binomial). The SEs are all larger than those for the fit of the naive Equation 13, but more similar to the SEs for the expanded LP model of Equation 16 (model B). The test of treatment effect is now F = 4.80, similar to that obtained with Equation 16 and much smaller than the F for the naive Equation 13 (model A). The use of quasi-likelihood for model fitting is very straightforward and can be applied to many modeling situations. It is often considered a simple fix for overdispersion in GLMs or GLMMs ( McCullagh and Nelder, 1989). Nevertheless, (Stroup 2013, pp. 347–348) was somewhat critical of the approach because an actual likelihood function is not used in the estimation as there is no true conditional distribution. That is, there is no model for the direct stochastic generation of the data. There are alternative approaches, such as the use of Equation 16 (model B), which may be consistent with the experimental design (i.e., random block and plot effects in this example). Model B (Equation 16) can be thought of as characterizing one possible data generation process for the given experimental design (e.g., the random effects of blocks and experimental units such as plots) and model C (Equation 17) as characterizing a consequence of the random effects affecting π (or μ in general) that are not in the model. Piepho (1999) and Madden et al. (2002) compared these and other GLMMs for the analysis of plant disease incidence data. Piepho (1999) should also be consulted for a more thorough discussion of data generation processes that would lead to model B or C. 6.5 Different distributions with overdispersion For unbounded count data, Poisson is commonly replaced by the NB to account for overdispersion at the experimental unit level (or for the clustering of individual observations, in general) (Madden et al., 2007). The conditional Poisson would be replaced with the conditional NB in Equation 13. Because the NB has its own overdispersion parameter k that is estimated, then a v[ij] term would not be added and a ϕ not specified for a quasi-distribution. This substitution of the NB for the Poisson is easy to do with the GLIMMIX procedure in SAS. For overdispersed binomial-type data (when there is an upper bound of n for y), such as in our example, the beta-binomial distribution (“Betabin”) can replace the binomial. This distribution is not in the exponential family and is much less used for mixed-model analysis, although it has had more use with GLMs with only fixed effects (Hughes and Madden, 1995; Morel and Neerchal, 2010) and for quantifying aspects of the spatial heterogeneity of plant diseases (Madden et al., 2007, 2018). To use this model, the ${y}_{ij}|{b}_{j} ~\mathit{\text{Bin}}\left({\pi }_{ij},n\right)$ in Equation 13 would be replaced with ${y}_{ij}|{b}_{j} ~\ mathit{\text{Betabin}}\left({\pi }_{ij},n, \varrho \right)$, where $\varrho$ is one formulation of the overdispersion parameter for this distribution (at $\varrho$ = 0, the Betabin reduces to the binomial; as $\varrho$ increases, overdispersion increases). The LP would contain only the treatment and block effect (as in Equation 13). This can be considered one type of DGLMM (Lee et al., 2006). The beta-binomial distribution is not available in the GLIMMIX procedure, but can be fitted using ML (not the REML) with the NLMIXED procedure when there are random effects (such as block effects). Considerably more coding is required to achieve this for the beta-binomial. An alternative to fitting the beta-binomial is to utilize the h-likelihood (see below). The results for the example with the beta-binomial as the conditional distribution (model D) are given in Table 3 based on ML estimation. Although the means are similar, the SEs are larger than those for model A (Equation 13) as the extra-binomial variability is taken into account, but smaller than those for models B and C (which take into account the extra-binomial variability in different ways). The smaller SEs for the beta-binomial are due, in part, to ML rather than REML being used in the estimation. Table 3 Table 3 Estimated means and standard errors (SEs, in parentheses) both on the linear predictor scale [logit(mean)] and the data scale (inverse link; with SEs based on the delta method) when generalized linear mixed models (GLMMs) were fitted to the number of diseased wheat spikes^a using maximum likelihood with the quadrature method for integral approximation for Equation 16 (model B-quad), for Equation 13 but with the beta-binomial for the conditional distribution of y instead of the binomial (model D), and when using Bayesian estimation with Equation 16 (model B-Bayes). 7 Fitting GLMMs to data Unlike with normality-based LMMs, the method for fitting GLMMs to data (i.e., parameter estimation and random effects prediction) is controversial (Stroup, 2013; Stroup and Claassen, 2020). Even the labels for the model-fitting methods are confusing and are used differently by different researchers, with the labels even evolving over time. For LMMs, REML is the generally preferred (and noncontroversial) method of model fitting, partly because it produces unbiased estimates (fixed effects and their SEs) or less biased estimates than ML (McCulloch and Searle, 2001; Galecki and Burzykowski, 2013; Stroup et al., 2018). For simple cases, REML duplicates the results for the mean square (and moment)-based methods that predate the contemporary likelihood-based methods for LMMs. Regular ML can also be used for LMMs. However, it is well known that this produces biased estimates; in particular, the variance estimates and the SEs will be too small for ML estimation, resulting in the test statistics being too large. The latter is an issue when the number of independent observations is small. REML and ML are both iterative methods, although simple situations can result in convergence in a single iteration. Recall that, for any experiment or survey, the observed data are a manifestation of the marginal distribution. Thus, the estimation requires integration over all the random effects in a conditional model to approximate the marginal distribution. However, there is no analytical solution (i.e., mathematical expression) for this integration with GLMMs; only approximations are possible, which is one of the complexities of working with this class of models. In fact, practical use of GLMMs only became possible with the availability of fast computers with large memory (and excellent programming). Despite the many labels in the literature, two broad frequentist methods can be used: model approximation (linearization) and integral approximation. There are also some more specialized methods. 7.1 Linearization Linearization involves a first-order Taylor series expansion of the inverse link function of the LP equation [π[ij] = g^−1(η[ij]) = h(η[ij])], generally centered on the current (or initial) estimates of the fixed and random effects (Breslow and Clayton, 1993; Wolfinger and O’Connell, 1993). A so-called pseudo-variable (y*) is formed based on this expansion, which is a function of the current parameter estimates (and random effect predictions) and the first derivative of the inverse link with respect to the LP. The expected value of y* is modeled as a function of the fixed and random effects (the right-hand side of the LP). The variance of y* conditional on the random effects is a complex function of the first derivatives and the variance function for the conditional distribution [e.g., nπ(1 − π) for binomial and μ for Poisson]. Wolfinger and O’Connell (1993) developed a pseudo-likelihood (PL) method based on linearization. The PL method makes the (approximating) assumption that the conditional distribution of y* (but not y ), given the random effects, is normal with a complicated variance (Stroup, 2013; Xie and Madden, 2014). As a consequence, the marginal distribution is also normal. Thus, one can use the machinery of normality-based LMMs to fit GLMMs and then (automatically) recover the relevant statistics for the actual distribution of y after convergence of the PL algorithm. The PL fitting algorithm is doubly iterative in that the fitting of the y* pseudo-variable is iterative (as is any LMM), and then the y* is updated at the end of each LMM fit based on the LMM results, with the LMM fitting being repeated with the updated y*. The double iterations continue until convergence (defined in different ways). With PL, one ultimately is basing the analysis on the actual distributions specified in the model for y (the normality assumptions for y* is only for the intermediate computational work). Restricted PL is the default in the GLIMMIX procedure in SAS (called RSPL in SAS), which uses REML-based model fitting for the pseudo-variable. As an option, the ML version of PL (ML-based PL; known as MSPL in GLIMMIX) can be performed. An overdispersion term, ϕ, can also be added as an option (see Equation 17) with the PL method (either RSPL or MSPL), which then becomes a quasi-likelihood approach (no longer a true statistical distribution for the conditional distribution of y). We are not aware of any R packages for PL estimation. Breslow and Clayton (1993) took a quasi-likelihood approach, which only assumes that the objective function to maximize in the iterative estimation process has the general form of a member of the exponential family of conditional distributions (defined through a mean and variance, but with a multiplicative overdispersion parameter). This approach is often labeled as penalized quasi-likelihood (PQL) or marginal quasi-likelihood (MQL). Implementation in R can be done with the “glmmPQL” function in the MASS package. In this program, the overdispersion parameter ϕ is always estimated (cannot be restricted to 1); therefore, the results do not correspond to a (conditional) distribution for y, but to a quasi-likelihood. The PL approach taken to fit model C in the example (Table 2) has analogy with the PQL method of Breslow and Clayton (1993). The PL method is extremely flexible, allowing not only true distributions but also quasi-likelihoods and can easily handle correlated observations, such as in temporal or spatial repeated measures ( Stroup, 2013). The RSPL method was therefore used to fit models A, B, and C to the example data discussed so far. 7.2 Integral approximation Instead of approximating the model to obtain a marginal distribution of a pseudo-variable, the integration over the random effects (of the original model) can be approximated to obtain an estimated marginal distribution of y. Generally, the best method to use is the Gauss–Hermit quadrature (quadrature for short), although it is computationally demanding and can be extremely (or painfully) slow for moderate to large data problems, sometimes requiring many hours if there are several factors and interactions. The Laplace approximation, a special case of quadrature, is another integral approximation method that works well for many datasets and often gives estimates very similar to those obtained with the more accurate quadrature (Joe, 2008; Stroup, 2013; Ruíz et al., 2023). Laplace can be very fast and works when quadrature is not possible or practical. Both Laplace and quadrature are available as options in GLIMMIX of SAS, and quadrature with one quadrature point (which reduces to one way of expressing the Laplace approximation) is the default in the lme4 package of R when fitting GLMMs. Linearization methods are not done with the “lme4” package. There are two important points worth noting with this approach: 1) Integral approximations are ML-based; that is, there is no restricted/residual (REML-like) version that reduces bias. The problems that come with the use of ML instead of REML with normality-based models carry over to GLMMs (bias of parameter estimates and SEs, test statistics, etc.). 2) Integral approximations require a true likelihood; therefore, quasi-likelihood-based models (e.g., ϕ > 1 with quasi-binomial likelihood) cannot be fitted. In particular, model C for overdispersion with the binomial data cannot be fitted using these approaches. 7.2.1 Example with integral approximation The results from model B (Equation 16, with a random v[ij] for plot effects) using quadrature are given in Table 3 (model B-quad). The results for the means are similar in this example to those obtained with RSPL (see Table 2), although the SEs are slightly smaller with the integral approximation methods (e.g., SEs of 0.245 versus 0.283 for the logit mean of treatment 1). The variance estimates are a little smaller here (0.175 versus 0.248 for ${\stackrel{^}{\sigma }}_{v}^{2}$), and the F statistic is a little larger for quadrature compared with that of RSPL. The smaller variances and SEs are a consequence of the use of ML- rather than REML-based methods. If the MSPL method in GLIMMIX was used, the results would be more similar between PL and integral approximations (both being ML-based). The linearization and integral approximation (quadrature) methods, however, will never give identical results. 7.3 Comparison of linearization and integral approximations Early on after the linearization methods were proposed, it was accepted as common knowledge that integral approximations are more accurate (less biased) than linearization when fitting GLMMs to data (Stroup, 2013). However, this conclusion was mostly based on assessments under extreme conditions [such as when there was only one individual (diseased or healthy) in an experimental unit]. The “common knowledge” has not held up based on more recent assessments, at least not as a generality (Couton and Stroup, 2013; Claassen, 2014; Piepho et al., 2018; Stroup and Claassen, 2020). Estimation performance has been recently assessed in detail by Stroup and Claassen (2020) (see also their online supplements) with extensive simulations, and the overall results showed that there is no overall best method for fitting GLMMs. The best method depends, in part, on how well the conditional distribution (such as binomial or Poisson) can be approximated by a normal distribution (see Figure 1). For conditional binomial data, this partly depends on n (the number of individual observations that gives a proportion y/n) for an individual experimental unit, such as a plot (block × treatment combination for a RCBD). When n is much less than 30, the conditional binomial may be quite skewed (especially when π is close to 0 or 1), and the linearization methods may give biased results. A small n per experimental unit coupled with a large number of experimental units (e.g., blocks in a RCBD) would be a situation where the likelihood approximation works best as the conditional distribution need not be normal-like. Integral approximations may also be desirable when π is very close to 0 or 1 at moderate values of n. On the other hand, the linearization methods, particularly the RSPL, are very accurate when n is 40 or higher in each experimental unit and may produce considerably less biased estimates of the means and SEs than the integral approximations, especially with the small number of experimental units (blocks) that are common in agricultural sciences. There are circumstances when all estimation methods are less than satisfactory (e.g., small n combined with a small number of replicates). Since integral approximations are strictly ML-based and not REML-based, they suffer from small-sample bias (small number of replicates, such as blocks in the RCBD case). Thus, with the typical number of four to six blocks in field experiments, integral approximations will typically lead to estimates of variances that are too small, leading to the SEs for means being too small and the test statistics for the significance of factor effects (overall or individual contrasts) being too large. These all result in a large type I error rate (rejecting the null hypothesis when the null hypothesis is true). The RSPL linearization method (with the same model) performs much better in these situations. The exception is when the n is very small, in which no method may be acceptable. However, the MSPL linearization method (or quasi-likelihood) would also suffer from the small-sample bias because it is ML- and not REML-based. A much more detailed investigation is given in Stroup and Claassen (2020). There are other advantages of the linearization methods. We have found when analyzing multiple datasets that the linearization methods more readily converge under a wide range of situations. Since these methods use normal distributions (as intermediate approximations), some methods strictly for LMMs are available. One is the use of the Kenward–Roger method to obtain more appropriate estimates of the SEs of the parameter estimates and adjustments for the denominator degrees of freedom for small samples (Stroup, 2013). Another advantage is that the algorithm does not “blow up” when one or more variance estimates are 0. Negative variances can even be obtained with the linearization methods as long as the variance of the marginal distribution is positive (Stroup et al., 2018). This allows for better control of type 1 errors in some circumstances. In contrast, with integral approximations, an estimate of 0 for a variance leads to nonsensical SEs (and other results). For the latter, the terms in the GLMM with a 0 variance must be removed and the model refitted. Based on Claassen (2014) and Stroup and Claassen (2020), one can make the following general recommendations for fitting GLMMs to conditional binomial data: ● For large n and small number of replicates (e.g., blocks): Use RSPL. ● For large n and large number of replicates: Use either approach as both have similar performance. ● For small n and small number of replicates: No method has shown great performance. RSPL may converge more easily. ● For small n and large number of replicates: Use integral approximation approaches. ● When using integral approximation, quadrature may be theoretically better than the Laplace approximation, but the Laplace approach gives similar results under many circumstances. Quadrature may not be computationally feasible in some situations, especially with several variances in the GLMM. ● If one wants to simply inflate the variance of the conditional distribution (ϕ > 1) to account for the overdispersion (extra-binomial variation due to not explicitly accounting for some sources of variation in the model), the binomial can be replaced with a quasi-likelihood (analogous to PL with the ϕ allowed to take any value). With a small number of replicates (blocks), this should be done with restricted PL (REML-based, such as RSPL in GLIMMIX). Similar recommendations can be made for conditional Poisson data (Stroup and Claassen, 2020). That is, if the number of replicates is small (as in many field studies), RSPL linearization outperforms, or strongly outperforms, the quadrature and Laplace integral approximation methods. The superior performance of RSPL may be less pronounced when the conditional distribution is highly skewed (typically when μ is very close to 0), but it still performs well. Integral approximations perform well primarily when there are many replicates. 7.4 Goodness of fit Even when linearization may be advantageous based on the bias of the parameter and variance estimates, among others, there are some other reasons to use integral approximations. In particular, the goodness of fit of models, including evaluation of the model assumptions, cannot easily or directly be done using linearization methods. This is because one is dealing with the (restricted) log-likelihood for the pseudo-variable (with RSPL or MSPL), not the (restricted) log-likelihood for the original observations. The goodness of fit of the model for the pseudo-variable, although done, is not a reliable metric to assess the GLMM (Littell et al., 2006). For instance, with the example in Table 2, adding the v[ij] block × treatment interaction to the GLMM (Equation 16) resulted in a point estimate of ${\sigma }_{v}^{2}$ greater than 0, with larger SEs for the means and smaller F statistic for significance compared with those of the model without this term (Equation 13). This suggests that the added random effect term was necessary. The results between the two models would be about the same (e.g., no increase in SEs) if the v[ij] term was not needed. However, this conclusion is only informal. The v[ij] term (or any random-effects terms) can be formally tested if the (restricted) log-likelihood for the data is available, which is obtained with the integral approximation. A likelihood-based confidence interval can be constructed for ${\sigma }_{v}^{2}$, and a chi-squared test of significance (null hypothesis of ${\sigma }_{v}^{2}$ = 0) is straightforward with the GLIMMIX procedure. A simpler approach with integral approximation approaches is to compare the AIC statistics for the model fits with and without the v[ij] term. With the example above, the AIC is 194.8 with the v[ij] term (model B-quad), while it is 229.7 with the simpler model [model A-quad (results not shown for A-quad)]. The lower AIC is an indication that there is a random plot effect that needs to be accounted for. Likelihood-based methods are especially appropriate for studying the magnitude of estimated variances when the focus is on random effects rather than on fixed effects. For instance, Kriss et al. (2012) used a hierarchical GLMM to investigate the variability of the incidence of Fusarium head blight among counties, fields within counties, and sampling sites within fields. Using a so-called small-area sampling approach, the authors demonstrated the use of likelihood-based confidence intervals for the variances in which the spatial scales had the greatest (and lowest) heterogeneity of disease. Presently, there are new computational methods for incorporating survey weights in the random effects of GLMMs to account for unequal sampling probabilities when analyzing survey data ( Diaz-Ramirez et al., 2020). Recent research by Piepho (2019, 2023) has shown how to estimate a coefficient of determination (R^2) for the fit of LMMs and GLMMs. An advantage of this new approach is that linearization or integral approximation methods can be used to fit the GLMM. Interested readers should consult these papers for a discussion of other proposals for R^2 calculations with mixed models. Using the algorithm in Piepho (2023), R^2 = 0.507 for the fit of Equation 16 (model B) to the example data. In our view, the calculation of relative measures of goodness of fit such as R^2 is most desirable when fitting LMMs or GLMMs with continuous explanatory variables, when trying to compare mixed models with different explanatory variables or number of explanatory variables, or when analyzing observational data rather than data from randomized trials. An example would be a random coefficient model for crop yield in relation to disease intensity measurement, with location–year as a random effect (Madden and Paul, 2009). In principle, graphical methods can be used to assess the goodness of fit, as well as select the most reasonable model, as they commonly are for normality-based LMMs (Littell et al., 2006; Stroup, 2013; Stroup et al., 2018). This includes assessment of different types of residual plots. There are many possible choices in viewing residuals for GLMMs, such as those determined on the model (i.e., linear predictor) or the data (i.e., inverse link) scale, or the type of scaling (standardization) to use for the residuals (e.g., raw, Pearson, or studentized). A useful presentation on this can be found in Stroup et al. (2018) and Stroup (2013). We agree with the authors of the former that more research is needed to fully evaluate the fit of GLMMs, especially when trying to decide on the most appropriate model to use. 7.5 The h-likelihood As an alternative to the linearization and integral methods described above, the so-called h-likelihood method is also possible (Lee and Nelder, 1996; Lee et al., 2006). Instead of estimating a marginal likelihood (either on a pseudo-variable scale or the original data scale), a joint likelihood is defined for both fixed and random effects. This approach is especially useful when one wants to specify that the random effects have non-normal distributions. An example would be a beta-distribution for the random block effect. Thus, the h-likelihood method is well suited for DGLMMs, such as when the beta-binomial distribution is used; however, in principle, h-likelihood can be applied to a wide range of problems. Although there is an R package (“hglm”), the methodology has still not been broadly adopted in statistics for routine data analysis with GLMMs, and some view the approach with skepticism (Meng, 2009). For normally distributed random effects, we have found that the h -likelihood estimation performs very similarly to RSPL in terms of fixed-effects parameter estimates, SEs, and other statistics when the data have a conditional binomial distribution (Piepho et al., 7.6 Bayesian estimation Although we emphasize frequentist methods in this paper, a Bayesian approach can be taken to fit a GLMM if one has an assumed true conditional distribution (binomial or Poisson, for instance) (Piepho and Madden, 2022). Thus, the Bayesian approach follows directly from the integral approximation methods in the sense that the model is not approximated and true distributions are used. Bayesian analysis is a different philosophy, but is becoming quite popular in many fields. Prior distributions (indicating the degree of certainty or uncertainty before the experiment) on all the parameters (e.g., τ[i]), random effects (e.g., b[j]), and the variance and covariance parameters (e.g., ${\sigma }_{b}^{2}$) need to be specified in order to ultimately estimate the marginal posterior distributions and credible intervals for the parameters and the differences of parameters (such as the means and the differences of means). Most Bayesian approaches (but not all) require a considerable amount of coding to implement, and is computationally expensive, where sampling of distributions is done, usually with the Markov chain Monte Carlo (MCMC) algorithm, until convergence is reached (although it is not always reached). The approach is much more challenging for a non-statistician. There are a range of R packages for Bayesian analysis (rstan), and the release of the BGLIMM procedure in SAS greatly facilitates Bayesian analysis with GLMMs as the syntax is very close to that of GLIMMIX. For those interested, Section 2.3 of Brown and Prescott (2015) is a good place to start, at least in terms of mixed models. The online tutorial by Stroup (2021) is extremely informative for the fitting of GLMMs using the BGLIMM procedure. Piepho and Madden (2022) further demonstrated the use of this Bayesian procedure for conducting a meta-analysis of conditional binomial data. 7.6.1 Example When fitting Equation 16 to the example data (model B-Bayes) with non-informative prior distributions, the treatment means were similar to those found with the frequentist analysis (Table 3). As is typical, the means were not the same, however. The standard deviations of the posterior distributions (analogous to the SEs in the frequentist analysis) were larger than those for the fit of the same model with quadrature (compare models B, B-quad, and B-Bayes). This is expected because Bayesian analysis takes into account the uncertainty of the variance estimates. 8 Overview of some additional GLMMs and general guidance for analysis 8.1 Marginal model Although one would normally be interested in the conditional means when fitting a GLMM to discrete data (Stroup, 2013), such as π[i], there may be situations when marginal means are desired. This type of GLMM can be fitted by expanding on the concepts explained above for marginal LMMs (Equations 7-11). The approach uses the quasi-likelihood and expands on the common method used with GLMs (no random effects) for repeated measures (Liang and Zeger, 1986), generally called generalized estimating equations (GEEs). As with LMMs, the LP for a RCBD consists only of the fixed effects, $\text{ logit}\left({\pi }_{ij}\right)={\eta }_{ij}=\theta +{\tau }_{i}$. The effects of block and plot are accounted for through a so-called working correlation or covariance matrix for the vector of the response variable for the j-th block, y[j] (the response vector would have six elements with six treatments). The working covariance matrix is not a true covariance matrix, but behaves like one. Details are given in Stroup (2015) and Stroup (2013). Fitting a marginal GLMM (with RSPL) to the example disease incidence data would produce the same mean proportions (after applying the inverse link function to the estimated logits), as found in the first column of results in Table 1 for a normality-based LMM fitted directly to the proportion data (Stroup, 2013). This is not unexpected because the marginal means and conditional means are the same with LMMs, as discussed above. The SEs of the estimated means with marginal GLMM do vary with the mean, as required, due to the dependency of variances on the means for non-normal data. Although similar here, the means from a marginal GLMM generally will not agree with the means from a conditional GLMM. Since we feel that researchers are mostly interested in conditional means (see above discussion), we do not give any more details on the marginal approach. 8.2 GLMMs for two expansions of the RCBD It is instructive to see how the GLMM for a RCBD is expanded for two different scenarios. 8.2.1 Sampling With field studies, it is common to collect multiple samples within each experimental unit. For instance, Madden et al. (2002) analyzed several datasets with a GLMM for the effects of fungicide treatment on the incidence of Phomopsis leaf blight in strawberry. Within each experimental unit (plot), there were either three or five samples, each consisting of n = 15 leaflets. Although the values from the clusters could be pooled to obtain a single y and n for each plot, sampling within plots can also be explicitly accounted for. The LMM linear predictor would be: $\begin{array}{ll} {\eta }_{ijk}=\text{θ}+{\tau }_{i}+{b}_{j}+{v}_{ij}& \text{(18A)}\end{array}$ where μ[ijk] = η[ijk], in which the ijk subscript represents the i-th treatment, the j-th block, and the k-th sampling unit within the plot (block × treatment combination), respectively, and v[ij] is the random plot effect (equivalent to the block × treatment interaction). The residual for the LMM (i.e., the variance of the conditional normal distribution, ${\stackrel{^}{\sigma }}_{e}^{2}$) represents the effect of the k-th sampling unit within the ij-th plot (block × treatment × sampling unit interaction). For a GLMM with a conditional binomial distribution, the simple use of Equation 18A, where logit(π[ijk]) = η[ijk], would fail to account for the effects of the sampling units within plots. This approach would be a generalization of model A (Equation 13). A quasi-likelihood approach could be taken using Equation 18A with the overdispersion parameter ϕ, producing a generalization of model C ( Equation 17; quasi-conditional binomial likelihood). Alternatively, the conditional model could be expanded to: $\begin{array}{ll}\text{ }{\eta }_{ijk}=\text{θ}+{\tau }_{i}+{b}_{j}+{v}_{ij}+{w}_{ijk}& \text{(18B)}\end{array}$ where w[ijk] is the random effect of the k-th sampling unit within the ij-th plot (analogous to the residual in a LMM, now a three-way interaction of block, treatment, and sampling unit). Now, one has a conditional binomial for y[ijk]: Both v[ij] and w[ijk] are assumed to have normal distributions with variances ${\sigma }_{v}^{2}$ and ${\sigma }_{w}^{2}$, respectively, when using Equation 18B. This would be an expansion of model B (Equation 16). Based on previous work by Piepho (1999), Madden et al. (2002) compared these and some other more complex GLMMs for the analysis of five strawberry datasets for disease incidence. Additional levels in the sampling hierarchy for disease incidence were analyzed for another crop/disease system by Piepho (1999) using data initially analyzed by Hughes and Madden (1995) using GLMs (all fixed effects). Only the RSPL (without or with a ϕ overdispersion term) was used by Piepho (1999) and Madden et al. (2002). Recall from above that the use of the ϕ parameter with PL is actually a form of quasi-likelihood. 8.2.2 Split plots A split plot with blocking is a very common design for field experiments in agriculture, involving two or more fixed-effects factors (Steel and Torrie, 1960; Schabenberger and Pierce, 2002). Here, one has three sizes of experimental units: blocks, whole plots (one of the fixed-effects factors, where the factor levels are randomized within each block), and sub-plots (another fixed-effects factor, where the factor levels are randomized within each whole plot). For instance, Moraes et al. (2022) studied the effect of cultivar resistance level and fungicide treatment on disease incidence for Fusarium head blight, as well as for other response variables assumed to have distributions in the exponential family. For a single year, cultivar was the whole plot (large experimental units), randomized within blocks, and fungicide treatment was the sub-plot (small experimental units), randomized within each whole plot factor level within blocks. With random blocks and a split-plot design, the LP for a normality-based LMM is: $\begin{array}{ll}\text{ }{\eta }_{ijk}=\text{θ}+{\gamma }_{i}+{\tau }_{j}+{\left(\gamma \tau \right)}_{ij}+{b}_{k}+{v}_{ik}& \text{(19A)}\end{array}$ where μ[ijk] = η[ijk], in which γ[i] is the effect of the i-th whole-plot factor, τ[j] is the effect of the j-th sub-plot factor, (γτ)[ij] is the interaction of the whole-plot and sub-plot factors, b [k] is the random effect of the k-th block, and v[ik] is the random effect of the ik-th whole plot (equivalent to the interaction of block and the whole-plot effect). The variance of the conditional normal distribution represents the variability among sub-plots within whole plots (equivalent to the interaction of block, whole-plot, and sub-plot effects). For a GLMM for disease incidence (conditional binomial), with logit(π[ijk]) = η[ijk], direct use of Equation 19A [an expansion of model A (Equation 13) for RCBD] would not account for the variability of the sub-plots. The extra-binomial variability due to a sub-plot effect could be accounted for by using quasi-likelihood for the binomial, i.e., using Equation 19B, but with a quasi-binomial likelihood with a ϕ parameter [expansion of model C (Equation 17)]. Alternatively, reflecting the random effects of sub-plots within whole plots and blocks, the LP can be expanded to: $\begin{array}{ll} {\eta }_{ijk}=\text{θ}+{\gamma }_{i}+{\tau }_{j}+{\left(\gamma \tau \right)}_{ij}+{b}_{k}+{v}_{ik}+{w}_{ijk}& \text{(19B)}\end{array}$ where w[ijk] is the random effect of the j-th sub-plot within the ik-th whole plot (equivalent to the interaction of block, whole plot, and sub-plot). This would be an expansion of model B (Equation 16) for a conditional binomial. Both v[ij] and w[ijk] are assumed to have normal distributions with variances ${\sigma }_{v}^{2}$ and ${\sigma }_{w}^{2}$, respectively. Equation 19B is the basis for the model used in Moraes et al. (2022), although the authors expanded it to account for an analysis of the split plot over years (where year and interactions with year were also random effects in the 8.3 Overall guidance for GLMMs with different designs GLMMs can be expanded in numerous ways to account for different treatment and experimental designs. Readers should consult Gbur et al. (2012); Stroup (2013); Stroup et al. (2018), and Ruíz et al. (2023) for examples. For those who are used to analyzing data with LMMs, here is the basic rule of thumb: ● Start with the LP equation that would be appropriate for a normality-based LMM analysis (there are many textbooks with this). ● For GLMMs involving either the binomial or Poisson conditional distribution, either: o Add a random-effects term to the LP that would be analogous to the residual term of a LMM, which can then be fitted with a linearization or an integral approximation method (see above for guidance) o Keep the LP of the LMM and use a quasi-likelihood (with the ϕ parameter) instead of a true conditional distribution. As described above, there are other approaches as well, such as the use of marginal models or the switch to other conditional distributions for overdispersed binomial-type data (e.g., beta-binomial). For repeated measures, or more generally when there is a variance–covariance matrix, the use of GLMM methods is more challenging, but can be done with a working covariance matrix (a GEE-type quasi-likelihood approach). See Stroup et al. (2018) and Stroup (2013) for discussion and examples. When the conditional distribution is not Poisson or binomial, such as the gamma or beta for (assumed) non-normal continuous random variables, there is actually a variance or scale parameter for the conditional distribution that is estimated from the data (Gbur et al., 2012). Thus, one might not need to consider adding a random effect term to the LP or use a quasi-likelihood approach. There may still be additional variability to account for, however, and readers should refer to Gbur et al. (2012). 8.4 The overall guidance applied to the example Several GLMMs and model-fitting methods were useful in analyzing the example dataset that exhibited the typical situation of higher variability than can be represented by the naive model A, the model that transfers directly from an LMM-based analysis. The use of model A would give a false sense of treatment effects with artificially low SEs. The addition of an experimental unit random-effects term (analogous to the residual in an LMM) was consistent with the experimental design and led to more realistic SEs and tests of significance. Due to the small number of blocks (typical for field studies in the plant sciences) and the relatively large number of observations within plots, the RSPL (model B), a linearization approach, was preferable to the quadrature (model B-quad), an integral approximation. Quasi-likelihood methods could also be used to account for the overdispersion exhibited by the data (model C), although one no longer has a true conditional statistical distribution in the model. A more complicated approach of changing the conditional distribution from binomial to beta-binomial was also effective in accounting for the overdispersion (model D), although this method will probably not be for routine analysis by non-statisticians. Except for model A, the differences in the results (i.e., the means and SEs) were not large, reflecting the fact that the estimated π values for the different treatments were not too close to 0 or 1. 9 Challenges 9.1 All zeros or all ones Although it remains popular to fit LMMs to transformed proportion data (for conditional binomial), care must be taken regarding the selected transformation. Transformations of observations such as logs or logits fail (i.e., are undefined) when the proportion is 0 or 1, and the log transformation fails when the proportion is 0 (Piepho, 2003) (the angular transformation is defined for these proportions). A major advantage of GLMMs for conditional binomial data is that the observed individual zeros or ones for proportions do not pose any difficulty; no ad hoc methods (adjustments for zeros and ones) are needed to accommodate the extreme values of the response variable. Nevertheless, there remains a potential problem with GLMMs fitted to data with a conditional binomial (or Poisson) distribution. For a RCBD and conditional binomial data, as an example, model fitting fails if all the y observations for a given treatment are equal to 0 or n (so that the proportions are 0 and 1 for all observations) (Gianinetti, 2020). Statistically, the problem is known as quasi-complete separation (Albert and Anderson, 1984; Clark et al., 2023). Since this is typically found with a very effective treatment for controlling disease, so that all y values are 0, this is sometimes called the all-zero problem, but it all applies to the situation where all the y values are equal to n. When one tries to use a likelihood-based model-fitting method with these data, the iterative procedure will never converge, or apparent convergence may be obtained, but for nonsensical or meaningless estimates of some means and absurdly large (and meaningless) values of some SEs (Albert and Anderson, 1984). Essentially, this is because with quasi-complete separation, mathematically, the ML estimates of one or more parameters are infinite (Heinze and Schemper, 2002). That is, some parameter estimates do not exist as real numbers when fitting a model using ML. Gianinetti (2020) and Claassen (2014) discussed the problem in more detail. One solution is to remove the treatments with all zeros (or all ns) and then refit the model. The problem is that the treatment resulting in all zeros often is of most interest, such as a new or a very effective pesticide. The treatment with all ns may be the control treatment that the investigator wishes to compare with others. The ad hoc approach is to add a small constant, c, to y (y′ = y + c) and add 2c to n (n′ = n + 2c). This gives values of prop′ (=y′/n′) that may be very close to 0 and 1, but never equal to these limits. The y′ data can then be analyzed. Of course, the choice of c will affect the results. A common choice for c is 0.5. Still, there are questions. In particular, should one modify all observations in this manner (all treatments and blocks) or just for the “problem” treatments? In fact, the modification is only needed for a single observation, and not all replicates, to avoid the problem (one of the blocks for the problem treatment). For GLMs (i.e., all fixed effects), models can be fitted to data of this type using a penalized likelihood method (Firth, 1993; Heinze and Schemper, 2002). For a RCBD, this would mean treating block as a fixed effect. This is straightforward with the LOGISTIC procedure in SAS and with some R packages. Although this may be reasonable for this balanced simple case, expansions to split plots and other designs (such as repeated measures and cluster sampling, among others) are not possible as random effects are required to account for aspects of the experimental designs. It turns out that Firth’s penalized likelihood method is related to the ad hoc approach of fitting a model to the modified y′ data (Firth, 1993). Firth (1993) also showed the link between his penalized likelihood approach and Bayesian analysis with certain types of prior distributions. Claassen (2014) has made important contributions by expanding on the approach in Firth (1993) for GLMMs for some situations. As far as we are aware, this generalization is not available in standard GLMM software. Overall, there is no best approach to recommend at this point for data with a conditional binomial (or Poisson) distribution, and a lot more research is needed. For the practicing data analyst, the ad hoc y′/n′ approach will be the easiest to implement in the GLMM context, by far. Alternatively, analyzing the proportion data (suitably transformed) using a normality-based LMM might be the most practical when there are multiple treatments with all zeros or all ns. The situation is different for continuous non-normal variables in the exponential family. For instance, with the commonly used gamma distribution (Gbur et al., 2012), the response variable is a non-zero positive value (0 is not allowed). With the beta distribution for continuous proportions, the response variable is between 0 and 1 (where 0 and 1 are not allowed). When using these as the conditional distribution in a GLMM, any individual observations outside the permitted range become missing values (just like any 0 becomes a missing value when logs are used in a LMM). An adjusted y (y′) would need to be used to keep the response variable within the required range, although the form of the adjustment will affect the results. There are also generalizations of the gamma and beta distributions that allow for zeros (or zeros and ones) (e.g., Basak and Balakrishnan, 2012), but these are outside the usual realm of GLMMs with exponential family distributions. More research is needed to deal with this problem when using GLMMs. 9.2 Too many zeros Misspecification of the conditional distribution is one of the causes of overdispersion for discrete data. Sometimes, there too many zeros in the dataset compared with what is possible to represent with a binomial distribution (or a beta-binomial). There could also be too many observations with all n individuals with disease (or whatever the trait of interest). For plant disease in the field, suppose, as a simple example, that there is no pathogen inoculum present in some plots (e.g., certain sections of a field), but there is inoculum in the other plots (e.g., fungal spores in soil); the presence or absence of inoculum is unknown to the investigator. This can be considered a consequence of the spatial heterogeneity of the unknown inoculum levels. Where inoculum is present, the conditional binomial may be appropriate for any given treatment and block (where y could be from 0 to n), but without the inoculum, there can be no infection (where y must be 0). For the conditional distribution, one then has a mixture of two distributions or two processes. One distribution would be binomial and the other would be a degenerate discrete distribution, where Prob(y = 0) = 1. This is known as the zero-inflated process, which is usually manifested by a relatively large frequency of zeros in a histogram compared with the rest of the frequency distribution at a given π (Cohen, 1960; Lambert, 1992). In Equation 13, as an example, all sub-equations are the same, except that we write ${y}_{ij}|{b}_{j} ~ \text{Mixture}\left({\pi }_{ij},n;\text{ω}\right)$, where “Mixture” is an arbitrary notation for the mixture distribution and ω is a so-called mixture probability parameter. The estimate of ω represents the proportion of the mixture that is of the binomial type, and 1 − ω represents the proportion that is the type with all zeros. Detection of zero inflation will typically be done only when there are many observations, such as when a dataset consists of multiple cluster samples within each experimental unit (plot). Failure to account for this situation will lead to biased estimates of the π[i] probabilities for treatments. There would typically be insufficient data in a simple RCBD, with one observation per combination of treatment and block, to identify and model zero inflation as a mixture distribution. There are excellent methods for fitting zero-inflated models with GLMs (all fixed effects), such as in GENMOD and FMM in SAS. It is much more challenging to fit zero-inflated GLMMs. In SAS, one would need to utilize nonlinear mixed-model software and write the code in NLMIXED. “glmmTMB” is an R package for zero-inflated discrete data, while “NBZIMM” is a mixed-model package for zero-inflated NB data (recall that NB is a generalization of the Poisson for overdispersed count data). Bayesian approaches may be very beneficial (Zuur et al., 2012). Readers should be on the lookout for new procedures and packages for fitting these types of GLMMs. 10 Example GLMMs in the plant sciences in agriculture We have assembled a list of papers where GLMMs for non-normal data (or data assumed to be non-normal) have been used in plant pathology, agronomy, and horticulture (Table 4). We also included papers when normal (Gaussian) data were analyzed together with non-normal data (for different response variables). This list is based on a search in Google Scholar and Web of Science. The list is only intended to give a sense of the range of applications of GLMMs and related data analyses and is not meant to be an exhaustive compilation. The response variables analyzed in these experiments can be placed in the following categories: i) continuous symmetric or asymmetric data (e.g., amount of mycotoxins, insect body length, seed weight, sucrose concentration, yield, disease severity, leaf damage, or beetle survival time); ii) discrete proportion data (e.g., number of leaves infected, germinated seeds, pigmented flowers, or hatched eggs out of a total of n individuals); and iii) count (discrete) data without a definable upper limit (e.g., number of insects, aphids, female eggs, or tubers). These studies were conducted to assess either the efficacy of fungicides or insecticides on disease and insect control or the effect of environmental or host factors on disease severity, mycotoxin accumulation, plant density, seed germination, seed weight, etc. A handful of studies are focused on a theoretical evaluation of the different elements of GLMMs using existing datasets (e.g., Piepho, 2019). Table 4 Table 4 Summary of some applications of generalized linear mixed models in the analysis of data from designed experiments in agricultural and horticultural systems. Binomial, NB, gamma, and Poisson were the most commonly used conditional distributions in these studies, while beta (Bilka et al., 2021; Susi and Laine, 2021) and Bernoulli (a special case of binomial with n = 1) (De Silva et al., 2014) were less frequently used. The link functions specified in the models in these studies were logit, log, probit, complementary log–log, and identity (i.e., no transformation of the mean). While several papers provided details on how model fitting was conducted and specified the estimation method [ML (quadrature or Laplace), PL, or quasi-likelihood)] used, other papers did not provide details of the estimation method used. We made no attempt to establish whether the model, the link function, or the estimation method specified in these studies was appropriately specified by the authors to match the data and the aims of individual studies. Furthermore, we did not evaluate whether the authors interpreted the results appropriately. This summary shows that the application of GLMMs in the above disciplines is diverse and is used to address hypotheses in a wide range of experimental conditions. We expect this to grow as researchers appreciate the clear argument of matching the model to the data (for a given experimental and treatment design) when using GLMMs for analysis. 11 Summary and conclusions There are several reasons to use GLMMs for the analysis of non-normal data, some of which were discussed in this paper, with many more details in several references (Littell et al., 2006; Bolker et al., 2009; Zuur et al., 2009; Gbur et al., 2012; Stroup, 2013; Brown and Prescott, 2015; Stroup et al., 2018; Gianinetti, 2020; Li et al., 2023; Ruíz et al., 2015, 2023). Perhaps the biggest argument is that one is better off matching the model to the data with a GLMM rather than changing (i.e., transforming) the data to match the model, as with a LMM (Gbur et al., 2012; Stroup, 2013). With some contemporary software, it is deceptively easy now to switch from normality-based LMMs to non-normality-based GLMMs. For instance, this simple model statement in the GLIMMIX procedure of SAS can be used to fit a LMM to the response variable y for a RCBD (as in our example): $\text{model} y=\text{ trt};$ (Random effects, such as blocks, are given in separate statements, and the identification of explanatory variables as factors is done in a separate statement). Using the same procedure, a conditional–binomial GLMM with a logit link function can be fitted to the same data with only one minor change to the model statement: $\text{model\hspace{0.17em}}y/n=\text{ trt};$ with all other statements (not shown here) being the same. The procedure automatically identifies this as a conditional binomial with logit link when the left-hand side of the equation is given as the ratio of a response variable to the number of individuals (i.e., y/n). Of course, the distribution (dist=) and link (link=) can be explicitly specified as options with this procedure (Gbur et al., 2012; Brown and Prescott, 2015; Ruiz et al., 2023), allowing for many different conditional distributions and links. However, this simple switch would be fitting the naive model A (Equation 13) to the RCBD data (assuming the block random effect was also specified). This is because, with an LMM, a residual is always estimated to account for the plot variability (in a field study) after accounting for other sources of variability, but is not estimated with a conditional binomial (or conditional Poisson). The variance is always nπ(1 − π) for a conditional binomial and always μ for the conditional Poisson. In separate statements for random effects, the block × treatment random term (v[ij]) would need to be added to fit model B (Equation 16), or the conditional distribution to be changed from binomial to a quasi-binomial likelihood with the ϕ overdispersion parameter to account for overdispersion (model C; Equation 17). Alternatively, a more complex conditional distribution could be used, such as the beta-binomial (but not with the GLIMMIX procedure). The important point here is that one should never lose sight of the experimental design when analyzing data (Piepho et al., 2003; Bello et al., 2004, 2016), as well as how to account for the relevant aspects of the design (fixed effects, random effects, and correlations) in the model (e.g., right-hand side of the equation for η, plus distributions for random effects) for any selected conditional distribution, whether using LMMs or GLMMs. Readers should consult Stroup (2013, 2015) for a more thorough discussion on this topic. Adding terms to a GLMM (or modifying conditional distributions) may be needed to account for the sources of variability that are automatically accounted for with LMMs. In this paper, we emphasized the switch from LMM to GLMM, with the assumption that the reader is more familiar with LMMs. We mostly restricted our attention to one type of response variable [discrete observations (counts) that can be represented as proportions or percentages] and one experimental design, namely, a RCBD. This approach was chosen to show the similarities and differences between LMMs and GLMMs, focusing on one small dataset. The differences in the results between LMMs and GLMMs will be greater in other datasets with larger random effect variances and when some of the π[i] values are closer to 0 or 1 for some treatments (Stroup, 2013). The differences between LMMs and GLMMs carry through to more complex experimental and treatment designs (Gbur et al., 2012; Ruíz et al., 2023). Issues in data analysis that are mostly well resolved for LMMs remain debatable with GLMMs. For instance, appropriate model fitting (estimation) can be controversial for GLMMs (Stroup and Claassen, 2020). In fact, some models, such as those involving a quasi-likelihood (e.g., models with an overdispersion parameter, ϕ), can only be fitted using certain estimation methods (either pseudo-likelihood or quasi-likelihood). For other models, investigators have choices for model fitting, and the best choice depends on the circumstances, as presented above. Despite earlier criticisms, linearization methods, particularly the restricted versions such as RSPL (analogous to REML for LMMs), have been shown recently to often be the best choice for model fitting, especially for discrete data and the small sample sizes typically used in field studies in agriculture (Stroup and Claassen, 2020). Integral approximations of the marginal likelihood (such as quadrature and Laplace) do have advantages when there are reasonably large numbers of replicates, especially when one is interested in comparing the goodness of fit of different models through differences in the log-likelihood or have specific interest in the variances and covariances. A key distinction of LMMs and GLMMs is the notion of conditional versus marginal means (Stroup, 2013). Conditional and marginal means are the same for normality-based LMMs, but are different, in general, for GLMMs. The magnitude of the random effect variances has a large influence on the magnitude of the difference between the conditional and marginal means. We agree with Stroup (2013, 2015) and Stroup et al. (2018) that most researchers would find the conditional means as being more meaningful and of primary interest, thus emphasizing the importance of conditional models in a GLMM-based analysis. The use of a conditional model, however, still allows for broad or narrow inference when testing hypotheses or calculating the confidence intervals for parameters or parameter differences. Inference space (broad vs. narrow) has to do with whether inference is over the entire population of random effects (say, all possible blocks, locations, etc.) or only for the specific random effect levels in the study. Littell et al. (2006) and Stroup (2013) should be consulted for more detail on this topic. Author contributions LM: Conceptualization, Formal analysis, Writing – original draft, Writing – review & editing. PO: Conceptualization, Data curation, Writing – review & editing. The author(s) declare that no financial support was received for the research, authorship, and/or publication of this article. Conflict of interest The authors declare that the research was conducted in the absence of any commercial or financial relationships that could be construed as a potential conflict of interest. The author(s) declared that they were an editorial board member of Frontiers, at the time of submission. This had no impact on the peer review process and the final decision. Publisher’s note All claims expressed in this article are solely those of the authors and do not necessarily represent those of their affiliated organizations, or those of the publisher, the editors and the reviewers. Any product that may be evaluated in this article, or claim that may be made by its manufacturer, is not guaranteed or endorsed by the publisher. Supplementary material The Supplementary Material for this article can be found online at: https://www.frontiersin.org/articles/10.3389/fhort.2024.1423462/full#supplementary-material Adkins S., McCollum T. G., Albano J. P., Kousik C. S., Baker C. A., Webster C. 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Keywords: binomial, conditional distribution, fixed effects, integral approximation, marginal distribution, pseudo-likelihood, random effects, quasi-likelihood Citation: Madden LV and Ojiambo PS (2024) The value of generalized linear mixed models for data analysis in the plant sciences. Front. Hortic. 3:1423462. doi: 10.3389/fhort.2024.1423462 Received: 25 April 2024; Accepted: 04 June 2024; Published: 25 June 2024. Edited by: Xiangming Xu , National Institute of Agricultural Botany (NIAB), United Kingdom Copyright © 2024 Madden and Ojiambo. This is an open-access article distributed under the terms of the Creative Commons Attribution License (CC BY). The use, distribution or reproduction in other forums is permitted, provided the original author(s) and the copyright owner(s) are credited and that the original publication in this journal is cited, in accordance with accepted academic practice. No use, distribution or reproduction is permitted which does not comply with these terms. *Correspondence: Laurence V. Madden, madden.1@osu.edu
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How To Do Long Division? - Steps and Methods - Chimpvine How To Do Long Division? – Steps and Methods Understanding Division Division is a fundamental operation in mathematics that involves the distribution of a quantity into equal parts. It is the process of splitting a number into equal groups or determining how many times one number is contained within another. In this article, we will explore the concept of division, division properties, and the method of long division. The Division Equation The division equation is represented as a ÷ b = c, where ‘a’ is the dividend, ‘b’ is the divisor, and ‘c’ is the quotient. The dividend is the number being divided, the divisor is the number by which the dividend is divided, and the quotient is the result of the division. When the number isn’t excatly divisor, there may be a remainder as well. It is the number that cannit be further How to Do Long Division Division is one of the four basic mathematical operations which includes addition, subtraction, and multiplication. Long division is a method used to divide large numbers or polynomials by hand. The process involves several steps, including dividing, multiplying, subtracting, and bringing down digits. It is a systematic approach to division that allows for the accurate calculation of the quotient and remainder. Steps of Doing Long Division Step 1: Set up the division problem. Step 2: Divide the first digit of the dividend by the divisor and write the answer on the top which will be the quotient. Step 3: Subtract the result from the number and write down the difference, Step 4: Bring down the next digit of the dividend Step 5: Repeat the same process until it can’t be furter divided. Let’s look at some examples. Example 1: Long Division without remainder Dividing 650 by 5 Step 1: Setting up the division problem Step 2: Divide the first digit 6, by 5 . As 6 cannot be divided by 5, we will look for a number less than 6 that 5 can divide, which is 5 itself. So, 5 ÷ 5 = 1, we will write 1 in the quotient. Step 3: Subtract 6 – 5 = 1. Step 4: We bring down the next digit of the dividen wihch is 5. The new dividend will be 15, we can divide 15 by 5 with quotient 3, as 5 * 3 = 15. Step 5: Subtract 15 – 15 =0, we will bring the next digit down, which is 0. Step 6: As 5 ÷ 0 = 0, we will now divide the dividend of 0 by 0, which will give us 0, so our answer will be 130. Example 2: Long Division with Remainder Divide 263 by 3 Step 1: Setting up the division problem Step 2: Divide the first digit 2, by 3 . As 3 is greater than 2, we will put 0 and bring the next digit down. Step 3: Then, 26 will be the dividend. We can’t divide 26 by 3, however we know that 3*8 = 24, so we will use 8 as the quotient. Step 4: Subtract 26 – 24 = 2. Bring down 3, so the next dividend will be 23. Step 5: We can’t divide 23 by 3, but 7 * 3 = 21, so we will use 7 as the quotient. Step 6: Subtract 23 – 21 = 2, there is no othher digit, so 2 is our remainder and we cannot go further. Learn division with interactive activities on our site ChimpVine. 1. Long Division for Two Digits Tip: If the divisor is a two digit number, we will look for the first two digit of the dividend and carry out the division. 2. Dividing with 0 Tip: If 0 is divided by any number, the answer will be 0. For example, 0 ÷55 = 0 3. Dividing with 1 Tip: If any number is divided by 1, the quotient will be the same number. For example, 40 ÷1=40 4. Division Formula Tip: We can use the division formula to verify quotient and remainder. The formula is Dividend = (Divisor * Quotient) + Remainder. Story: “The Long Division Adventures of Sam and Emma” Sam and Emma were two curious explorers who encountered real-life situations where the concept of long division played a crucial role in solving problems and making decisions. Challenge 1: The Recipe Dilemma Sam and Emma were passionate about cooking and decided to bake a batch of cookies. The recipe called for 2 cups of flour, and they needed to divide the flour equally into 4 portions. By using the division equation, they determined that each portion should contain 0.5 cups of flour, ensuring the perfect consistency for their cookies. Challenge 2: The Budgeting Puzzle As they planned a road trip, Sam and Emma needed to divide their total budget of $500 over 5 days of travel. By applying the division equation, they calculated that they could spend $100 per day, allowing them to manage their expenses effectively and enjoy their journey without overspending. Challenge 3: The Party Planning Adventure Sam and Emma had to divide 60 party favors into 6 equal groups for their friend’s birthday party. Using the concept of division, they determined that each group would receive 10 party favors, ensuring that all the guests had an equal share of the fun. The properties of division includes division by 1 Property, division by itself property, division by 0 property, and division of 0 by any number property. Division and multiplication are inverse operations, meaning they undo each other. When dividing a number by another, it is equivalent to multiplying the first number by the reciprocal of the second number. This relationship allows for the conversion of division problems into multiplication problems and vice versa. Long division is significant as it provides a systematic method for dividing large numbers or polynomials. It allows for the accurate calculation of the quotient and remainder, making it a valuable tool for solving complex division problems by hand. Yes, division is commonly used to solve real-life problems related to sharing, distribution, measurement, and resource allocation. It helps in determining equal portions, calculating rates, and managing quantities, making it an essential skill for everyday tasks and decision-making. Division performs operations with fractions and decimals, such as finding the quotient of two fractions or dividing a decimal by a whole number. It allows for the comparison, simplification, and conversion of fractions and decimals, making it a fundamental aspect of working with rational numbers. Like? Share it with your friends
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The Science Behind Relationship PredictorsThe Science Behind Relationship Predictors The Science Behind Relationship Predictors Relationship Predictor Survey There are now several tools that predict which genealogical relationships are most likely for given amounts of shared DNA. With the recent introduction of the cM Explainer^TM at MyHeritage, Blaine Bettinger and I launched a citizen science initiative to evaluate them. Volunteers have been contributing match data for known relationships then running each match through each of the tools in the study. The goal is to determine which predictor, if any, is best. You can read more about the project here. To date, the survey has more than 6,000 entries for 50 unique relationships ranging from parent–child to 8th cousins twice removed. We’re putting the finishing touches on the analyses and are excited to share them with the genealogy community soon. In the meantime, it helps to understand how relationship predictors work. What Could Go Wrong? Biology likes to throw curveballs, so even the best DNA-based predictions will be off sometimes. On average, 1st cousins (1Cs) share roughly 875 cM, but a 1st cousin once removed (1C1R) can share that much on the rare occasion. If you happen to be that rare 1C1R, all of the known predictive tools will peg you as a 1C over a 1C1R. That doesn’t mean the tools are faulty—nor that you are!—just that this is a particularly tough case. Of course, not all predictive tools are equal. They are based on computer models, and those models make assumptions. If the assumptions are wrong, the predictions can be off as well. The Assumptions By understanding those assumptions, we can get a better sense of how much credence to give a predictive tool. Let’s consider some of them. Genome Size For simplicity, in this post I will assume that a parent–child match shares 3500 cM and that the amount is halved each generation. Sex-specific Inheritance Patterns We inherit exactly 50% (3500 cM) of our autosomal DNA from each parent but not exactly 25% (1750 cM) from each grandparent. That’s because a process called “crossing over” in a parent’s body divvies up the grandparents’ DNA before passing it on, and that division is rarely equitable. Although the average is ≈1750 cM, some grandparent–grandchild matches will share more and some will share less. Distributions of atDNA shared with maternal (pink) and paternal (blue) grandparents. The higher the bar, the more likely that centimorgan amount is for the relationship. It turns out that egg production involves more crossover events than sperm production. Just as you’re more likely to get a 50-50 split of heads versus tails if you flip a coin 20 times than 10 times, you’re more likely to share close to the average of 1750 cM with your maternal grandparents than your paternal ones; the crossover events are analogous to coin flips. That means relationship prediction for your maternal relatives should be slightly more accurate, because those matches will be closer to the average, on average. Family Structure How likely a DNA match is to be any given relationship depends, in part, on how many such relatives you have. For example, a match of 1750 cM could be a grandparent, grandchild, aunt or uncle, niece or nephew, or half sibling, and which is more likely will depend on how many of each you have. On paper, I have four grandparents, one aunt, one uncle, and one half sister. Based on those numbers, a 1750-cM match to me has a 57% chance of being a grandparent, a 29% chance of being an aunt or uncle, and a 14% chance of being a half sibling. A younger cousin of mine has four grandparents, six aunts/uncles, and two half siblings, so for her the probabilities are 33%, 50%, and 17%, respectively. Those probabilities are quite different. Any tool that tries to differentiate those relationships must make assumptions about the family structure, usually based on what’s typical for the population. If your family structure doesn’t fit those assumptions, the predictions can be off for you yet be quite good for someone else. Population Growth Population growth also plays a key role here. Consider an unknown match who shares somewhere between the averages for 1C and 1C1R. Which of those two relationships is more likely depends, in part, on how many 1Cs and 1C1Rs you have. If each of your 1Cs had one child each, then you should have an equal number of 1Cs and 1C1Rs. (Let’s ignore our parents’ first cousins for now, just to get the point across.) The distributions will look like this: The break-even point, where a match is equally likely to be either relationship, is around 625 cM. However, if each of those family members had two kids each, you have twice as many 1C1Rs as 1Cs, and the distributions would look like this: In this case, the break-even point is closer to 650 cM and a 625-cM match is more likely to be a 1C1R. If your family averages four kids per couple, the probabilities are shifted even more. The most popular predictive tools are based on proprietary data from either AncestryDNA or MyHeritage, neither of which has publicly stated the population growth factor they use. However, it’s safe to assume it’s around 2.5, which is a fairly standard fertility rate for developed countries over the past century. Again, if your family doesn’t align well with that assumption, the predictors may not work as well for you. Perhaps the biggest challenge for relationship predictors is endogamy, the practice of marrying within a cultural or geographic group. Endogamy is common around the world and causes people to be related in more than one way. Those “extra” connections can increase the amount of shared DNA and throw off predictor tools that assume each match is related only once. Thus, those of us from endogamous populations will get less reliable relationship predictions. The Best Tool The ideal relationship predictor would be tailored to your particular family structure, your population’s growth rate, your level of endogamy, and so on. Such a tool does not exist yet. (But stay tuned!) In its absence, we would like to know which of the available tools gives most people the correct relationships most of the time. The next post in this series will reveal how the tools did! Learn More! Whether you love math or hate it, you’ll get more out of genetic genealogy if you understand some basic concepts and how they apply to our DNA. Join me for my upcoming class “No One Told Me There Would Be Math! DNA Numbers Made Easy.” It’s meant to be accessible to everyone, even if you haven’t done math since high school. The same lecture is offered on two dates/times, so sign up for the one that bet fits your schedule. Posts in This Series 8 thoughts on “The Science Behind Relationship Predictors” 1. Thank you for the elegant explanation. Family structure would seem to be the major factor making us misfits in the current models. I would love to contribute to research in that direction. 2. You say “your parent’s 1Cs had one child each” would count in the total number of 1C1R’s. However, the children of one of my parent’s 1st cousins are the same number of generations removed from my grandparents and are therefore my 2nd cousins, correct? The only way to get a 1C1R is the children of your 1st cousins, although technically, if we were to treat 1st cousins the same way we do other cousins, we could say we could refer to their parents as 1C1R as well, but we more commonly refer to them as aunts and uncles. 😉 1. Good catch! And a reminder that relationship prediction is often more complicated than a simplified description can capture. The children of your 1Cs are 1C1Rs, but so are your parent’s 1Cs. I’ll edit the text to be a bit more clear. 3. > The most popular predictive tools are based on proprietary data from either AncestryDNA or MyHeritage, neither of which has publicly stated the population growth factor they use. However, it’s safe to assume it’s around 2.5, which is a fairly standard fertility rate for developed countries over the past century. You should be able to calculate this from public trees, given estimated birthdates for currently-living people, and maybe given family, regional / cultural, and generational variations. 4. A small pedantic correction: I think you are more likely to get a(n exact) 50-50 split from 10 coin tosses than 20 (~17.6% against ~12.5%, if I’ve done the sums right). With 2 coin tosses its 50% of course. But we know what you mean: the distribution is “tighter” around 50-50 with increasing n. 1. Point taken! 5. Differentiating between the 3 bundled relationship categories within that 2nd (25%) degree of relatedness is what needs to be done to figure out how the testers are related. It can obviously be done successfully the way you describe either with advanced calculation methods or with one of the good 3rd party tools available. Another way to differentiate the relationship categories that belong to the 25% degree is by matching the generational separation between the two testers to the generational separation that the relationship categories describe. Grandparent Grandchild is a 2 generation separation category, Uncle-Aunt/Niece-Nephew is a 1 generation separation category and 1/2 sibling is a 0 generation separation category. In each degree of relatedness there won’t be more than one category per number of generations separated. The calculators rank probability based on the number of people related in that category, the more people the more likely the testers are to be related in that category which is a liklihood based on statistical probability. The alternate approach requires knowledge or an estimated guess at generational separation by the user to select the category that aligns to their situation with the other tester. I think the latter approach has a better chance of turning out to be correct on the first try just because it involves the user and makes them look at their specific situation in a way that no tool maker can anticipate. However as assisted reproduction becomes more and more prevalent (there goes your volume) our perception of generational separation based on age alone are less and less likely to be correct. I think at the moment the user has a better chance of picking the correct category by knowing or estimating how many generations separate the two testers than by looking at the chances based on the number of people possible in the category but our ability to guess at how many generations separate testers based on their ages is going to keep declining until, at some point in the future statistical probability based on an estimated head count overtakes the simplicity of eliminating the grandparent/grandchild category completely for two 20 year old testers and rolling the dice that the two 20 year old testers are 0 generation rather than 1 generation apart (because of course both are possible but 0 generation of separation would be more likely if we go back to estimating by a % of the population method). A hybrid of the two methodologies might work well now and in the future. The one thing that troubled me about the analogy of counting my own relatives in those categories to get a percentage of liklihood is that it does not take the fact that the other tester calculating from their side might have completely different percentages of liklihood which means counting known relatives is the wrong thing to be counting. The known relatives need to be excluded from the equation completely so that the focus is only on the unkown relationship vs those population factors you mention and of course the maximum path count (number of ways two people can be related in those categories). Anyway there are a lot of ways to go about the same problem. Thanks for explaining how most of the tools work! 1. If you have a priori reasons to gauge which generation(s) the matches are in, that’s a great way of narrowing the field. That can be hard to tell, though. For example, my grandfather was one of 14 children, with a spread of 25 years between the youngest and oldest. The spread in the next generation down is even bigger. An alternate way of going about it is to use multiple DNA matches to start to hone in on a generation. There’s no one-size approach. We need to use all the tools in our toolbox! This site uses Akismet to reduce spam. Learn how your comment data is processed.
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# programming mathematical methods for machine learning Words or phrase for the review: «programming mathematical methods for machine learning» COMPGI07: Programming and Mathematical Methods for Machine… » Browse Hierarchy COMPGI07: COMPGI07: Programming and Mathematical Methods for Machine Learning. Back to COMPS_ENG: Computer Science ... 5 Ways To Understand Machine Learning Algorithms Without a… » Aug 24, 2015… A lot of these tangential mathematical techniques are often bundled in… The need to being a better more capable programmer drives them to it. Machinelearningmastery.com Programming & Mathematical Methods for Machine Learning » COMPM012- Programming & Mathematical Methods for Machine Learning… complexity and linear algebra with an aim to its application to machine learning. Cs.ucl.ac.uk Mathematical Methods in Artificial Intelligence: Edward A. Bender… » Mathematical Methods in Artificial Intelligence introduces the student to the important… Paradigms of Artificial Intelligence Programming: Case Studies in Common Lisp… Math for Machine Learning: Open Doors to Data Science and Artificial ... Amazon.com CO-145 Mathematical Methods (Autumn 2015) | Statistical Machine… » Statistical Machine Learning… CO-145 Mathematical Methods (Autumn 2015)… in using basic mathematical techniques for later year courses in Computing. Wp.doc.ic.ac.uk Massimiliano Pontil - Courses » GI07/M012: Programming and Mathematical Methods for Machine Learning Fall 2014… GI01/4C55: Supervised Learning, Fall 2013, Dept of Computer Science. Www0.cs.ucl.ac.uk Mathematical Methods for Machine Learning - AMSI Summer School » Mathematical Methods for Machine Learning. Mathematical… Familiarity with a programming language for statistics, such as MATLAB, R or Python. Students ... Ss.amsi.org.au Math for Programmers - Manning » To score a job in data science, machine learning, computer graphics, and cryptography, you need to bring strong math skills to the party. Math for Programmers ... You're Doing it Wrong. Why Machine Learning Does Not Have to Be… » Mar 28, 2018… The frame of machine learning and even mathematics using the top-down… When I started to learn programming and software engineering, it was in…. These are methods used in machine learning for data projection, data ... Machinelearningmastery.com The Mathematics of Machine Learning – Towards Data Science » Mar 23, 2017… There are many reasons why the mathematics of Machine Learning is important and I… for understanding the optimization methods used for machine learning.… Coding the Matrix: Linear Algebra through Computer Science ... Towardsdatascience.com We may use cookies to offer you a better browsing experience, analyze site traffic, personalize content, and serve targeted advertisements. If you continue to use this site, you consent to our use of The information forward from this site may be provided by third parties. We will not be responsible with outside links, contents from source of information, methods of using, using or consequence of contents with users. All direct or indirect risk related to use of this site is borne entirely by you, the user. We use advertising companies as Google AdSense, to serve ads when you visit our website. These companies may use information (not including your name, address, email address, or telephone number) about your visits to this and other websites in order to provide advertisements about goods and services of interest to you. If you would like more information about this practice and to know your choices about not having this information used by these companies, see https://policies.google.com/technologies/ads.
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Combine VLOOKUP with SUMIF - Written by Puneet Yes, you can combine VLOOKUP and SUMIF. In SUMIF, there is a criteria argument in which you can use VLOOKUP to create a dynamic value. This will allow you to change the criteria by changing the lookup value in the VLOOKUP. The use of this combo formula is unique. In this tutorial, we will learn to combine SUMIF and VLOOKUP to create a formula. Combine SUMIF and VLOOKUP 1. First, in a cell enter “=SUMIF(“, for the range argument, refer to the product ID range that you have in table1. 2. After that, in the second augment, you need to use the VLOOKUP function to lookup for the product ID by using the product name from the cell above. 3. Next, in the third argument of the SUMIF, refer to the quantity column to use as a sum_range. 4. In the end, enter the closing parentheses and hit enter to get the result. How this Formula Works Let’s break this formula into three parts: In the first part, you have specified the range where you have the product ID. In the second part, you have the VLOOKUP that takes the product name from cell B15 and lookup it in the table2. For Headphones, we have the Product ID OT-356. In the third part, we have the quantity column as sum_range. In short, VLOOKUP helps you find the Product ID with the Product Name and then SUMIF takes that Product ID and looks it up in the Product ID column, and then sums values from the Quantity column. As I said, when you use SUMIF and VLOOKUP, this makes your formula a dynamic formula. When you change the product name in the cells it changes the result. SUMIF and VLOOKUP in Multiple Sheets You can use this combination even when you have both tables in multiple sheets differently. In the above example, you have the product name table on a different sheet. Leave a Comment
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