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custom_pipeline="speech_to_image_diffusion", |
speech_model=model, |
speech_processor=processor, |
torch_dtype=torch.float16, |
) |
diffuser_pipeline.enable_attention_slicing() |
diffuser_pipeline = diffuser_pipeline.to(device) |
output = diffuser_pipeline(speech_data) |
plt.imshow(output.images[0]) |
This example produces the following image: |
VQModel The VQ-VAE model was introduced in Neural Discrete Representation Learning by Aaron van den Oord, Oriol Vinyals and Koray Kavukcuoglu. The model is used in π€ Diffusers to decode latent representations into images. Unlike AutoencoderKL, the VQModel works in a quantized latent space. The abstract from the paper ... |
Tuple of downsample block types. up_block_types (Tuple[str], optional, defaults to ("UpDecoderBlock2D",)) β |
Tuple of upsample block types. block_out_channels (Tuple[int], optional, defaults to (64,)) β |
Tuple of block output channels. layers_per_block (int, optional, defaults to 1) β Number of layers per block. act_fn (str, optional, defaults to "silu") β The activation function to use. latent_channels (int, optional, defaults to 3) β Number of channels in the latent space. sample_size (int, optional, defaults... |
The component-wise standard deviation of the trained latent space computed using the first batch of the |
training set. This is used to scale the latent space to have unit variance when training the diffusion |
model. The latents are scaled with the formula z = z * scaling_factor before being passed to the |
diffusion model. When decoding, the latents are scaled back to the original scale with the formula: z = 1 / scaling_factor * z. For more details, refer to sections 4.3.2 and D.1 of the High-Resolution Image |
Synthesis with Latent Diffusion Models paper. norm_type (str, optional, defaults to "group") β |
Type of normalization layer to use. Can be one of "group" or "spatial". A VQ-VAE model for decoding latent representations. This model inherits from ModelMixin. Check the superclass documentation for itβs generic methods implemented |
for all models (such as downloading or saving). forward < source > ( sample: FloatTensor return_dict: bool = True ) β VQEncoderOutput or tuple Parameters sample (torch.FloatTensor) β Input sample. return_dict (bool, optional, defaults to True) β |
Whether or not to return a models.vq_model.VQEncoderOutput instead of a plain tuple. Returns |
VQEncoderOutput or tuple |
If return_dict is True, a VQEncoderOutput is returned, otherwise a plain tuple |
is returned. |
The VQModel forward method. VQEncoderOutput class diffusers.models.vq_model.VQEncoderOutput < source > ( latents: FloatTensor ) Parameters latents (torch.FloatTensor of shape (batch_size, num_channels, height, width)) β |
The encoded output sample from the last layer of the model. Output of VQModel encoding method. |
Paint by Example Paint by Example: Exemplar-based Image Editing with Diffusion Models is by Binxin Yang, Shuyang Gu, Bo Zhang, Ting Zhang, Xuejin Chen, Xiaoyan Sun, Dong Chen, Fang Wen. The abstract from the paper is: Language-guided image editing has achieved great success recently. In this paper, for the first time, ... |
Variational Auto-Encoder (VAE) model to encode and decode images to and from latent representations. image_encoder (PaintByExampleImageEncoder) β |
Encodes the example input image. The unet is conditioned on the example image instead of a text prompt. tokenizer (CLIPTokenizer) β |
A CLIPTokenizer to tokenize text. unet (UNet2DConditionModel) β |
A UNet2DConditionModel to denoise the encoded image latents. scheduler (SchedulerMixin) β |
A scheduler to be used in combination with unet to denoise the encoded image latents. Can be one of |
DDIMScheduler, LMSDiscreteScheduler, or PNDMScheduler. safety_checker (StableDiffusionSafetyChecker) β |
Classification module that estimates whether generated images could be considered offensive or harmful. |
Please refer to the model card for more details |
about a modelβs potential harms. feature_extractor (CLIPImageProcessor) β |
A CLIPImageProcessor to extract features from generated images; used as inputs to the safety_checker. π§ͺ This is an experimental feature! Pipeline for image-guided image inpainting using Stable Diffusion. This model inherits from DiffusionPipeline. Check the superclass documentation for the generic methods |
implemented for all pipelines (downloading, saving, running on a particular device, etc.). __call__ < source > ( example_image: Union image: Union mask_image: Union height: Optional = None width: Optional = None num_inference_steps: int = 50 guidance_scale: float = 5.0 negative_prompt: Union = None num_images_per_pro... |
An example image to guide image generation. image (torch.FloatTensor or PIL.Image.Image or List[PIL.Image.Image]) β |
Image or tensor representing an image batch to be inpainted (parts of the image are masked out with |
mask_image and repainted according to prompt). mask_image (torch.FloatTensor or PIL.Image.Image or List[PIL.Image.Image]) β |
Image or tensor representing an image batch to mask image. White pixels in the mask are repainted, |
while black pixels are preserved. If mask_image is a PIL image, it is converted to a single channel |
(luminance) before use. If itβs a tensor, it should contain one color channel (L) instead of 3, so the |
expected shape would be (B, H, W, 1). height (int, optional, defaults to self.unet.config.sample_size * self.vae_scale_factor) β |
The height in pixels of the generated image. width (int, optional, defaults to self.unet.config.sample_size * self.vae_scale_factor) β |
The width in pixels of the generated image. num_inference_steps (int, optional, defaults to 50) β |
The number of denoising steps. More denoising steps usually lead to a higher quality image at the |
expense of slower inference. guidance_scale (float, optional, defaults to 7.5) β |
A higher guidance scale value encourages the model to generate images closely linked to the text |
prompt at the expense of lower image quality. Guidance scale is enabled when guidance_scale > 1. negative_prompt (str or List[str], optional) β |
The prompt or prompts to guide what to not include in image generation. If not defined, you need to |
pass negative_prompt_embeds instead. Ignored when not using guidance (guidance_scale < 1). num_images_per_prompt (int, optional, defaults to 1) β |
The number of images to generate per prompt. eta (float, optional, defaults to 0.0) β |
Corresponds to parameter eta (Ξ·) from the DDIM paper. Only applies |
to the DDIMScheduler, and is ignored in other schedulers. generator (torch.Generator or List[torch.Generator], optional) β |
A torch.Generator to make |
generation deterministic. latents (torch.FloatTensor, optional) β |
Pre-generated noisy latents sampled from a Gaussian distribution, to be used as inputs for image |
generation. Can be used to tweak the same generation with different prompts. If not provided, a latents |
tensor is generated by sampling using the supplied random generator. output_type (str, optional, defaults to "pil") β |
The output format of the generated image. Choose between PIL.Image or np.array. return_dict (bool, optional, defaults to True) β |
Whether or not to return a StableDiffusionPipelineOutput instead of a |
plain tuple. callback (Callable, optional) β |
A function that calls every callback_steps steps during inference. The function is called with the |
following arguments: callback(step: int, timestep: int, latents: torch.FloatTensor). callback_steps (int, optional, defaults to 1) β |
The frequency at which the callback function is called. If not specified, the callback is called at |
every step. Returns |
StableDiffusionPipelineOutput or tuple |
If return_dict is True, StableDiffusionPipelineOutput is returned, |
otherwise a tuple is returned where the first element is a list with the generated images and the |
second element is a list of bools indicating whether the corresponding generated image contains |
βnot-safe-for-workβ (nsfw) content. |
The call function to the pipeline for generation. Example: Copied >>> import PIL |
>>> import requests |
>>> import torch |
>>> from io import BytesIO |
>>> from diffusers import PaintByExamplePipeline |
>>> def download_image(url): |
... response = requests.get(url) |
... return PIL.Image.open(BytesIO(response.content)).convert("RGB") |
>>> img_url = ( |
... "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/image/example_1.png" |
... ) |
>>> mask_url = ( |
... "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/mask/example_1.png" |
... ) |
>>> example_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/reference/example_1.jpg" |
>>> init_image = download_image(img_url).resize((512, 512)) |
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