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provided, text embeddings are generated from the prompt input argument. negative_prompt_embeds (torch.FloatTensor, optional) — |
Pre-generated negative text embeddings. Can be used to easily tweak text inputs (prompt weighting). If |
not provided, negative_prompt_embeds are generated from the negative_prompt input argument. return_dict (bool, optional, defaults to True) — |
Whether or not to return a AudioPipelineOutput instead of a plain tuple. callback (Callable, optional) — |
A function that calls every callback_steps steps during inference. The function is called with the |
following arguments: callback(step: int, timestep: int, latents: torch.FloatTensor). callback_steps (int, optional, defaults to 1) — |
The frequency at which the callback function is called. If not specified, the callback is called at |
every step. cross_attention_kwargs (dict, optional) — |
A kwargs dictionary that if specified is passed along to the AttentionProcessor as defined in |
self.processor. output_type (str, optional, defaults to "np") — |
The output format of the generated image. Choose between "np" to return a NumPy np.ndarray or |
"pt" to return a PyTorch torch.Tensor object. Returns |
AudioPipelineOutput or tuple |
If return_dict is True, AudioPipelineOutput is returned, otherwise a tuple is |
returned where the first element is a list with the generated audio. |
The call function to the pipeline for generation. Examples: Copied >>> from diffusers import AudioLDMPipeline |
>>> import torch |
>>> import scipy |
>>> repo_id = "cvssp/audioldm-s-full-v2" |
>>> pipe = AudioLDMPipeline.from_pretrained(repo_id, torch_dtype=torch.float16) |
>>> pipe = pipe.to("cuda") |
>>> prompt = "Techno music with a strong, upbeat tempo and high melodic riffs" |
>>> audio = pipe(prompt, num_inference_steps=10, audio_length_in_s=5.0).audios[0] |
>>> # save the audio sample as a .wav file |
>>> scipy.io.wavfile.write("techno.wav", rate=16000, data=audio) AudioPipelineOutput class diffusers.AudioPipelineOutput < source > ( audios: ndarray ) Parameters audios (np.ndarray) — |
List of denoised audio samples of a NumPy array of shape (batch_size, num_channels, sample_rate). Output class for audio pipelines. |
The Stable Diffusion Guide 🎨 |
Intro ... |
Stable Diffusion is a Latent Diffusion model developed by researchers from the Machine Vision and Learning group at LMU Munich, a.k.a CompVis. |
Model checkpoints were publicly released at the end of August 2022 by a collaboration of Stability AI, CompVis, and Runway with support from EleutherAI and LAION. For more information, you can check out the official blog post. ... |
Since its public release the community has done an incredible job at working together to make the stable diffusion checkpoints faster, more memory efficient, and more performant. ... |
🧨 Diffusers offers a simple API to run stable diffusion with all memory, computing, and quality improvements. |
This notebook walks you through the improvements one-by-one so you can best leverage StableDiffusionPipeline for inference. |
Prompt Engineering 🎨 |
When running *Stable Diffusion* in inference, we usually want to generate a certain type, or style of image and then improve upon it. Improving upon a previously generated image means running inference over and over again with a different prompt and potentially a different seed until we are happy with our generation. |
So to begin with, it is most important to speed up stable diffusion as much as possible to generate as many pictures as possible in a given amount of time. |
This can be done by both improving the computational efficiency (speed) and the memory efficiency (GPU RAM). |
Let’s start by looking into computational efficiency first. |
Throughout the notebook, we will focus on runwayml/stable-diffusion-v1-5: |
Copied |
model_id = "runwayml/stable-diffusion-v1-5" |
Let’s load the pipeline. |
Speed Optimization ... |
Copied |
from diffusers import StableDiffusionPipeline |
pipe = StableDiffusionPipeline.from_pretrained(model_id) |
We aim at generating a beautiful photograph of an old warrior chief and will later try to find the best prompt to generate such a photograph. For now, let’s keep the prompt simple: ... |
Copied |
prompt = "portrait photo of a old warrior chief" ... |
To begin with, we should make sure we run inference on GPU, so let’s move the pipeline to GPU, just like you would with any PyTorch module. |
Copied |
pipe = pipe.to("cuda") ... |
To generate an image, you should use the [~StableDiffusionPipeline.__call__] method. |
To make sure we can reproduce more or less the same image in every call, let’s make use of the generator. See the documentation on reproducibility here for more information. ... |
Copied |
generator = torch.Generator("cuda").manual_seed(0) |
Now, let’s take a spin on it. |
Copied |
image = pipe(prompt, generator=generator).images[0] |
image |
Cool, this now took roughly 30 seconds on a T4 GPU (you might see faster inference if your allocated GPU is better than a T4). |
The default run we did above used full float32 precision and ran the default number of inference steps (50). The easiest speed-ups come from switching to float16 (or half) precision and simply running fewer inference steps. Let’s load the model now in float16 instead. ... |
Copied |
import torch |
pipe = StableDiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16) |
pipe = pipe.to("cuda") |
And we can again call the pipeline to generate an image. |
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