website/docs/user-guide/skills/optional/mlops/mlops-stable-diffusion.md

---
title: "Stable Diffusion Image Generation"
sidebar_label: "Stable Diffusion Image Generation"
description: "State-of-the-art text-to-image generation with Stable Diffusion models via HuggingFace Diffusers"
---

{/* This page is auto-generated from the skill's SKILL.md by website/scripts/generate-skill-docs.py. Edit the source SKILL.md, not this page. */}

# Stable Diffusion Image Generation

State-of-the-art text-to-image generation with Stable Diffusion models via HuggingFace Diffusers. Use when generating images from text prompts, performing image-to-image translation, inpainting, or building custom diffusion pipelines.

## Skill metadata

| | |
|---|---|
| Source | Optional — install with `hermes skills install official/mlops/stable-diffusion` |
| Path | `optional-skills/mlops/stable-diffusion` |
| Version | `1.0.0` |
| Author | Orchestra Research |
| License | MIT |
| Dependencies | `diffusers>=0.30.0`, `transformers>=4.41.0`, `accelerate>=0.31.0`, `torch>=2.0.0` |
| Tags | `Image Generation`, `Stable Diffusion`, `Diffusers`, `Text-to-Image`, `Multimodal`, `Computer Vision` |

## Reference: full SKILL.md

:::info
The following is the complete skill definition that Hermes loads when this skill is triggered. This is what the agent sees as instructions when the skill is active.
:::

# Stable Diffusion Image Generation

Comprehensive guide to generating images with Stable Diffusion using the HuggingFace Diffusers library.

## When to use Stable Diffusion

**Use Stable Diffusion when:**
- Generating images from text descriptions
- Performing image-to-image translation (style transfer, enhancement)
- Inpainting (filling in masked regions)
- Outpainting (extending images beyond boundaries)
- Creating variations of existing images
- Building custom image generation workflows

**Key features:**
- **Text-to-Image**: Generate images from natural language prompts
- **Image-to-Image**: Transform existing images with text guidance
- **Inpainting**: Fill masked regions with context-aware content
- **ControlNet**: Add spatial conditioning (edges, poses, depth)
- **LoRA Support**: Efficient fine-tuning and style adaptation
- **Multiple Models**: SD 1.5, SDXL, SD 3.0, Flux support

**Use alternatives instead:**
- **DALL-E 3**: For API-based generation without GPU
- **Midjourney**: For artistic, stylized outputs
- **Imagen**: For Google Cloud integration
- **Leonardo.ai**: For web-based creative workflows

## Quick start

### Installation

```bash
pip install diffusers transformers accelerate torch
pip install xformers  # Optional: memory-efficient attention
```

### Basic text-to-image

```python
from diffusers import DiffusionPipeline
import torch

# Load pipeline (auto-detects model type)
pipe = DiffusionPipeline.from_pretrained(
    "stable-diffusion-v1-5/stable-diffusion-v1-5",
    torch_dtype=torch.float16
)
pipe.to("cuda")

# Generate image
image = pipe(
    "A serene mountain landscape at sunset, highly detailed",
    num_inference_steps=50,
    guidance_scale=7.5
).images[0]

image.save("output.png")
```

### Using SDXL (higher quality)

```python
from diffusers import AutoPipelineForText2Image
import torch

pipe = AutoPipelineForText2Image.from_pretrained(
    "stabilityai/stable-diffusion-xl-base-1.0",
    torch_dtype=torch.float16,
    variant="fp16"
)
pipe.to("cuda")

# Enable memory optimization
pipe.enable_model_cpu_offload()

image = pipe(
    prompt="A futuristic city with flying cars, cinematic lighting",
    height=1024,
    width=1024,
    num_inference_steps=30
).images[0]
```

## Architecture overview

### Three-pillar design

Diffusers is built around three core components:

```
Pipeline (orchestration)
├── Model (neural networks)
│   ├── UNet / Transformer (noise prediction)
│   ├── VAE (latent encoding/decoding)
│   └── Text Encoder (CLIP/T5)
└── Scheduler (denoising algorithm)
```

### Pipeline inference flow

```
Text Prompt → Text Encoder → Text Embeddings
                                    ↓
Random Noise → [Denoising Loop] ← Scheduler
                      ↓
               Predicted Noise
                      ↓
              VAE Decoder → Final Image
```

## Core concepts

### Pipelines

Pipelines orchestrate complete workflows:

| Pipeline | Purpose |
|----------|---------|
| `StableDiffusionPipeline` | Text-to-image (SD 1.x/2.x) |
| `StableDiffusionXLPipeline` | Text-to-image (SDXL) |
| `StableDiffusion3Pipeline` | Text-to-image (SD 3.0) |
| `FluxPipeline` | Text-to-image (Flux models) |
| `StableDiffusionImg2ImgPipeline` | Image-to-image |
| `StableDiffusionInpaintPipeline` | Inpainting |

### Schedulers

Schedulers control the denoising process:

| Scheduler | Steps | Quality | Use Case |
|-----------|-------|---------|----------|
| `EulerDiscreteScheduler` | 20-50 | Good | Default choice |
| `EulerAncestralDiscreteScheduler` | 20-50 | Good | More variation |
| `DPMSolverMultistepScheduler` | 15-25 | Excellent | Fast, high quality |
| `DDIMScheduler` | 50-100 | Good | Deterministic |
| `LCMScheduler` | 4-8 | Good | Very fast |
| `UniPCMultistepScheduler` | 15-25 | Excellent | Fast convergence |

### Swapping schedulers

```python
from diffusers import DPMSolverMultistepScheduler

# Swap for faster generation
pipe.scheduler = DPMSolverMultistepScheduler.from_config(
    pipe.scheduler.config
)

# Now generate with fewer steps
image = pipe(prompt, num_inference_steps=20).images[0]
```

## Generation parameters

### Key parameters

| Parameter | Default | Description |
|-----------|---------|-------------|
| `prompt` | Required | Text description of desired image |
| `negative_prompt` | None | What to avoid in the image |
| `num_inference_steps` | 50 | Denoising steps (more = better quality) |
| `guidance_scale` | 7.5 | Prompt adherence (7-12 typical) |
| `height`, `width` | 512/1024 | Output dimensions (multiples of 8) |
| `generator` | None | Torch generator for reproducibility |
| `num_images_per_prompt` | 1 | Batch size |

### Reproducible generation

```python
import torch

generator = torch.Generator(device="cuda").manual_seed(42)

image = pipe(
    prompt="A cat wearing a top hat",
    generator=generator,
    num_inference_steps=50
).images[0]
```

### Negative prompts

```python
image = pipe(
    prompt="Professional photo of a dog in a garden",
    negative_prompt="blurry, low quality, distorted, ugly, bad anatomy",
    guidance_scale=7.5
).images[0]
```

## Image-to-image

Transform existing images with text guidance:

```python
from diffusers import AutoPipelineForImage2Image
from PIL import Image

pipe = AutoPipelineForImage2Image.from_pretrained(
    "stable-diffusion-v1-5/stable-diffusion-v1-5",
    torch_dtype=torch.float16
).to("cuda")

init_image = Image.open("input.jpg").resize((512, 512))

image = pipe(
    prompt="A watercolor painting of the scene",
    image=init_image,
    strength=0.75,  # How much to transform (0-1)
    num_inference_steps=50
).images[0]
```

## Inpainting

Fill masked regions:

```python
from diffusers import AutoPipelineForInpainting
from PIL import Image

pipe = AutoPipelineForInpainting.from_pretrained(
    "runwayml/stable-diffusion-inpainting",
    torch_dtype=torch.float16
).to("cuda")

image = Image.open("photo.jpg")
mask = Image.open("mask.png")  # White = inpaint region

result = pipe(
    prompt="A red car parked on the street",
    image=image,
    mask_image=mask,
    num_inference_steps=50
).images[0]
```

## ControlNet

Add spatial conditioning for precise control:

```python
from diffusers import StableDiffusionControlNetPipeline, ControlNetModel
import torch

# Load ControlNet for edge conditioning
controlnet = ControlNetModel.from_pretrained(
    "lllyasviel/control_v11p_sd15_canny",
    torch_dtype=torch.float16
)

pipe = StableDiffusionControlNetPipeline.from_pretrained(
    "stable-diffusion-v1-5/stable-diffusion-v1-5",
    controlnet=controlnet,
    torch_dtype=torch.float16
).to("cuda")

# Use Canny edge image as control
control_image = get_canny_image(input_image)

image = pipe(
    prompt="A beautiful house in the style of Van Gogh",
    image=control_image,
    num_inference_steps=30
).images[0]
```

### Available ControlNets

| ControlNet | Input Type | Use Case |
|------------|------------|----------|
| `canny` | Edge maps | Preserve structure |
| `openpose` | Pose skeletons | Human poses |
| `depth` | Depth maps | 3D-aware generation |
| `normal` | Normal maps | Surface details |
| `mlsd` | Line segments | Architectural lines |
| `scribble` | Rough sketches | Sketch-to-image |

## LoRA adapters

Load fine-tuned style adapters:

```python
from diffusers import DiffusionPipeline

pipe = DiffusionPipeline.from_pretrained(
    "stable-diffusion-v1-5/stable-diffusion-v1-5",
    torch_dtype=torch.float16
).to("cuda")

# Load LoRA weights
pipe.load_lora_weights("path/to/lora", weight_name="style.safetensors")

# Generate with LoRA style
image = pipe("A portrait in the trained style").images[0]

# Adjust LoRA strength
pipe.fuse_lora(lora_scale=0.8)

# Unload LoRA
pipe.unload_lora_weights()
```

### Multiple LoRAs

```python
# Load multiple LoRAs
pipe.load_lora_weights("lora1", adapter_name="style")
pipe.load_lora_weights("lora2", adapter_name="character")

# Set weights for each
pipe.set_adapters(["style", "character"], adapter_weights=[0.7, 0.5])

image = pipe("A portrait").images[0]
```

## Memory optimization

### Enable CPU offloading

```python
# Model CPU offload - moves models to CPU when not in use
pipe.enable_model_cpu_offload()

# Sequential CPU offload - more aggressive, slower
pipe.enable_sequential_cpu_offload()
```

### Attention slicing

```python
# Reduce memory by computing attention in chunks
pipe.enable_attention_slicing()

# Or specific chunk size
pipe.enable_attention_slicing("max")
```

### xFormers memory-efficient attention

```python
# Requires xformers package
pipe.enable_xformers_memory_efficient_attention()
```

### VAE slicing for large images

```python
# Decode latents in tiles for large images
pipe.enable_vae_slicing()
pipe.enable_vae_tiling()
```

## Model variants

### Loading different precisions

```python
# FP16 (recommended for GPU)
pipe = DiffusionPipeline.from_pretrained(
    "model-id",
    torch_dtype=torch.float16,
    variant="fp16"
)

# BF16 (better precision, requires Ampere+ GPU)
pipe = DiffusionPipeline.from_pretrained(
    "model-id",
    torch_dtype=torch.bfloat16
)
```

### Loading specific components

```python
from diffusers import UNet2DConditionModel, AutoencoderKL

# Load custom VAE
vae = AutoencoderKL.from_pretrained("stabilityai/sd-vae-ft-mse")

# Use with pipeline
pipe = DiffusionPipeline.from_pretrained(
    "stable-diffusion-v1-5/stable-diffusion-v1-5",
    vae=vae,
    torch_dtype=torch.float16
)
```

## Batch generation

Generate multiple images efficiently:

```python
# Multiple prompts
prompts = [
    "A cat playing piano",
    "A dog reading a book",
    "A bird painting a picture"
]

images = pipe(prompts, num_inference_steps=30).images

# Multiple images per prompt
images = pipe(
    "A beautiful sunset",
    num_images_per_prompt=4,
    num_inference_steps=30
).images
```

## Common workflows

### Workflow 1: High-quality generation

```python
from diffusers import StableDiffusionXLPipeline, DPMSolverMultistepScheduler
import torch

# 1. Load SDXL with optimizations
pipe = StableDiffusionXLPipeline.from_pretrained(
    "stabilityai/stable-diffusion-xl-base-1.0",
    torch_dtype=torch.float16,
    variant="fp16"
)
pipe.to("cuda")
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
pipe.enable_model_cpu_offload()

# 2. Generate with quality settings
image = pipe(
    prompt="A majestic lion in the savanna, golden hour lighting, 8k, detailed fur",
    negative_prompt="blurry, low quality, cartoon, anime, sketch",
    num_inference_steps=30,
    guidance_scale=7.5,
    height=1024,
    width=1024
).images[0]
```

### Workflow 2: Fast prototyping

```python
from diffusers import AutoPipelineForText2Image, LCMScheduler
import torch

# Use LCM for 4-8 step generation
pipe = AutoPipelineForText2Image.from_pretrained(
    "stabilityai/stable-diffusion-xl-base-1.0",
    torch_dtype=torch.float16
).to("cuda")

# Load LCM LoRA for fast generation
pipe.load_lora_weights("latent-consistency/lcm-lora-sdxl")
pipe.scheduler = LCMScheduler.from_config(pipe.scheduler.config)
pipe.fuse_lora()

# Generate in ~1 second
image = pipe(
    "A beautiful landscape",
    num_inference_steps=4,
    guidance_scale=1.0
).images[0]
```

## Common issues

**CUDA out of memory:**
```python
# Enable memory optimizations
pipe.enable_model_cpu_offload()
pipe.enable_attention_slicing()
pipe.enable_vae_slicing()

# Or use lower precision
pipe = DiffusionPipeline.from_pretrained(model_id, torch_dtype=torch.float16)
```

**Black/noise images:**
```python
# Check VAE configuration
# Use safety checker bypass if needed
pipe.safety_checker = None

# Ensure proper dtype consistency
pipe = pipe.to(dtype=torch.float16)
```

**Slow generation:**
```python
# Use faster scheduler
from diffusers import DPMSolverMultistepScheduler
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)

# Reduce steps
image = pipe(prompt, num_inference_steps=20).images[0]
```

## References

- **[Advanced Usage](https://github.com/NousResearch/hermes-agent/blob/main/optional-skills/mlops/stable-diffusion/references/advanced-usage.md)** - Custom pipelines, fine-tuning, deployment
- **[Troubleshooting](https://github.com/NousResearch/hermes-agent/blob/main/optional-skills/mlops/stable-diffusion/references/troubleshooting.md)** - Common issues and solutions

## Resources

- **Documentation**: https://huggingface.co/docs/diffusers
- **Repository**: https://github.com/huggingface/diffusers
- **Model Hub**: https://huggingface.co/models?library=diffusers
- **Discord**: https://discord.gg/diffusers