| | --- |
| | license: cc-by-nc-4.0 |
| | library_name: diffusers |
| | base_model: PixArt-alpha/PixArt-XL-2-1024-MS |
| | tags: |
| | - lora |
| | - text-to-image |
| | inference: False |
| | --- |
| | # ⚡ Flash Diffusion: FlashPixart ⚡ |
| |
|
| |
|
| | Flash Diffusion is a diffusion distillation method proposed in [Flash Diffusion: Accelerating Any Conditional |
| | Diffusion Model for Few Steps Image Generation](http://arxiv.org/abs/2406.02347) *by Clément Chadebec, Onur Tasar, Eyal Benaroche, and Benjamin Aubin* from Jasper Research. |
| | This model is a **66.5M** LoRA distilled version of [Pixart-α](https://huggingface.co/PixArt-alpha/PixArt-XL-2-1024-MS) model that is able to generate 1024x1024 images in **4 steps**. |
| | See our [live demo](https://huggingface.co/spaces/jasperai/FlashPixart) and official [Github repo](https://github.com/gojasper/flash-diffusion). |
| |
|
| |
|
| | <p align="center"> |
| | <img style="width:700px;" src="assets/flash_pixart.jpg"> |
| | </p> |
| |
|
| | # How to use? |
| |
|
| | The model can be used using the `PixArtAlphaPipeline` from `diffusers` library directly. It can allow reducing the number of required sampling steps to **4 steps**. |
| |
|
| | ```python |
| | import torch |
| | from diffusers import PixArtAlphaPipeline, Transformer2DModel, LCMScheduler |
| | from peft import PeftModel |
| | |
| | # Load LoRA |
| | transformer = Transformer2DModel.from_pretrained( |
| | "PixArt-alpha/PixArt-XL-2-1024-MS", |
| | subfolder="transformer", |
| | torch_dtype=torch.float16 |
| | ) |
| | transformer = PeftModel.from_pretrained( |
| | transformer, |
| | "jasperai/flash-pixart" |
| | ) |
| | |
| | # Pipeline |
| | pipe = PixArtAlphaPipeline.from_pretrained( |
| | "PixArt-alpha/PixArt-XL-2-1024-MS", |
| | transformer=transformer, |
| | torch_dtype=torch.float16 |
| | ) |
| | |
| | # Scheduler |
| | pipe.scheduler = LCMScheduler.from_pretrained( |
| | "PixArt-alpha/PixArt-XL-2-1024-MS", |
| | subfolder="scheduler", |
| | timestep_spacing="trailing", |
| | ) |
| | |
| | pipe.to("cuda") |
| | |
| | prompt = "A raccoon reading a book in a lush forest." |
| | |
| | image = pipe(prompt, num_inference_steps=4, guidance_scale=0).images[0] |
| | ``` |
| | <p align="center"> |
| | <img style="width:400px;" src="assets/raccoon.png"> |
| | </p> |
| |
|
| | # Training Details |
| | The model was trained for 40k iterations on 4 H100 GPUs (representing approximately 188 hours of training). Please refer to the [paper](http://arxiv.org/abs/2406.02347) for further parameters details. |
| |
|
| | **Metrics on COCO 2014 validation (Table 4)** |
| | - FID-10k: 29.30 (4 NFE) |
| | - CLIP Score: 0.303 (4 NFE) |
| |
|
| | ## Citation |
| | If you find this work useful or use it in your research, please consider citing us |
| |
|
| | ```bibtex |
| | @misc{chadebec2024flash, |
| | title={Flash Diffusion: Accelerating Any Conditional Diffusion Model for Few Steps Image Generation}, |
| | author={Clement Chadebec and Onur Tasar and Eyal Benaroche and Benjamin Aubin}, |
| | year={2024}, |
| | eprint={2406.02347}, |
| | archivePrefix={arXiv}, |
| | primaryClass={cs.CV} |
| | } |
| | ``` |
| |
|
| | ## License |
| | This model is released under the the Creative Commons BY-NC license. |