Text-to-Image
Diffusers
Trained with AutoTrain
stable-diffusion-xl
stable-diffusion-xl-diffusers
lora
template:sd-lora
import torch
from diffusers import DiffusionPipeline
# switch to "mps" for apple devices
pipe = DiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-xl-base-1.0", dtype=torch.bfloat16, device_map="cuda")
pipe.load_lora_weights("kpal002/Dreambooth-SDXL")
prompt = "A portrait in a distinctive expressive style characterized by vivid color use, emotional depth, and dynamic brush strokes, capturing the nuanced expressions of children and teenagers."
image = pipe(prompt).images[0]AutoTrain SDXL LoRA DreamBooth - kpal002/Dreambooth-SDXL
Model description
These are kpal002/Dreambooth-SDXL LoRA adaption weights for stabilityai/stable-diffusion-xl-base-1.0.
The weights were trained using DreamBooth.
LoRA for the text encoder was enabled: False.
Special VAE used for training: None.
Trigger words
You should use A portrait in a distinctive expressive style characterized by vivid color use, emotional depth, and dynamic brush strokes, capturing the nuanced expressions of children and teenagers. to trigger the image generation.
Download model
Weights for this model are available in Safetensors format.
Download them in the Files & versions tab.
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Model tree for kpal002/Dreambooth-SDXL
Base model
stabilityai/stable-diffusion-xl-base-1.0