text stringlengths 60 353k | source stringclasses 2 values |
|---|---|
**Krull–Akizuki theorem**
Krull–Akizuki theorem:
In commutative algebra, the Krull–Akizuki theorem states the following: Let A be a one-dimensional reduced noetherian ring, K its total ring of fractions. Suppose L is a finite extension of K. If A⊂B⊂L and B is reduced, then B is a one-dimensional noetherian ring. Furthermore, for every nonzero ideal I of B, B/I is finite over A.Note that the theorem does not say that B is finite over A. The theorem does not extend to higher dimension. One important consequence of the theorem is that the integral closure of a Dedekind domain A in a finite extension of the field of fractions of A is again a Dedekind domain. This consequence does generalize to a higher dimension: the Mori–Nagata theorem states that the integral closure of a noetherian domain is a Krull domain.
Proof:
First observe that A⊂B⊂KB and KB is a finite extension of K, so we may assume without loss of generality that L=KB Then L=Kx1+⋯+Kxn for some x1,…,xn∈B Since each xi is integral over K, there exists ai∈A such that aixi is integral over A.
Let C=A[a1x1,…,anxn] Then C is a one-dimensional noetherian ring, and C⊂B⊂Q(C) , where Q(C) denotes the total ring of fractions of C.
Proof:
Thus we can substitute C for A and reduce to the case L=K Let pi be minimal prime ideals of A; there are finitely many of them. Let Ki be the field of fractions of A/pi and Ii the kernel of the natural map B→K→Ki . Then we have: A/pi⊂B/Ii⊂Ki and K≃∏Ki .Now, if the theorem holds when A is a domain, then this implies that B is a one-dimensional noetherian domain since each B/Ii is and since B≃∏B/Ii . Hence, we reduced the proof to the case A is a domain. Let 0≠I⊂B be an ideal and let a be a nonzero element in the nonzero ideal I∩A . Set In=anB∩A+aA . Since A/aA is a zero-dim noetherian ring; thus, artinian, there is an l such that In=Il for all n≥l . We claim alB⊂al+1B+A.
Proof:
Since it suffices to establish the inclusion locally, we may assume A is a local ring with the maximal ideal m . Let x be a nonzero element in B. Then, since A is noetherian, there is an n such that mn+1⊂x−1A and so an+1x∈an+1B∩A⊂In+2 . Thus, anx∈an+1B∩A+A.
Now, assume n is a minimum integer such that n≥l and the last inclusion holds. If n>l , then we easily see that anx∈In+1 . But then the above inclusion holds for n−1 , contradiction. Hence, we have n=l and this establishes the claim. It now follows: B/aB≃alB/al+1B⊂(al+1B+A)/al+1B≃A/(al+1B∩A).
Hence, B/aB has finite length as A-module. In particular, the image of I there is finitely generated and so I is finitely generated. The above shows that B/aB has dimension zero and so B has dimension one. Finally, the exact sequence B/aB→B/I→(0) of A-modules shows that B/I is finite over A. ◻ | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Molecular recognition feature**
Molecular recognition feature:
Molecular recognition features (MoRFs) are small (10-70 residues) intrinsically disordered regions in proteins that undergo a disorder-to-order transition upon binding to their partners. MoRFs are implicated in protein-protein interactions, which serve as the initial step in molecular recognition. MoRFs are disordered prior to binding to their partners, whereas they form a common 3D structure after interacting with their partners. As MoRF regions tend to resemble disordered proteins with some characteristics of ordered proteins, they can be classified as existing in an extended semi-disordered state.
Categorization:
MoRFs can be separated in 4 categories according to the shape they form once bound to their partners.The categories are: α-MoRFs (when they form alpha-helixes) β-MoRFs (when they form beta-sheets) irregular-MoRFs (when they don't form any shape) complex-MoRFs (combination of the above categories)
MoRF predictors:
Determining protein structures experimentally is a very time-consuming and expensive process. Therefore, recent years have seen a focus on computational methods for predicting protein structure and structural characteristics. Some aspects of protein structure, such as secondary structure and intrinsic disorder, have benefited greatly from applications of deep learning on an abundance of annotated data. However, computational prediction of MoRF regions remains a challenging task due to the limited availability of annotated data and the rarity of the MoRF class itself. Most current methods have been trained and benchmarked on the sets released by the authors of MoRFPred in 2012, as well as another set released by the authors of MoRFChibi based on experimentally-annotated MoRF data. The table below, adapted from, details some methods currently available for MoRF prediction (as well as related problems).
Databases:
mpMoRFsDBMutual Folding Induced by Binding (MFIB) database | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Train speed optimization**
Train speed optimization:
Train speed optimization, also known as Zuglaufoptimierung, is a system that reduces the need for trains to brake and accelerate, resulting in smoother and more efficient operation.
Train speed optimization:
While train speed optimization needs some technical infrastructure, it is more of an operational concept than a technical installation. One can relatively easily implement train speed optimization using for instance cab signalling (e.g. using ETCS), but the presence of a cab signalling system does not necessarily mean that it uses train speed optimization. Train speed optimization may also be implemented using conventional signalling.
Conventional signal operation:
Usually, trains are allowed to run at the maximum speed the track allows until the distant signal of next occupied block. This is inefficient in many cases, because this way the train comes to a halt in front of the red signal and has to accelerate again from zero.
Advantages using train speed optimization:
If the train slows down much earlier, given the right timing, it reaches the distant signal just when the home signal switches to green, and so does not need to stop. Thus, wear on the brakes is reduced and the train uses less energy. But the main reason, especially for trains that accelerate slowly, is that the train passes the home signal at high speed, compared to the conventional case where the train often has to accelerate from standstill. This effectively increases track capacity, because the time it takes for the train to run from the distant signal (that has just turned green) to the home signal is often much less than the time it takes for a train to accelerate from the home signal.
Equipment:
For a train speed optimization system to work, it is necessary to have a signalling system which is capable of displaying several different speeds, for instance 40, 60, 90 km/h and the full line speed, which also requires a train protection system that is able to handle these cases (cab signalling may replace these installations). Further, the track must be equipped with inductive loops that detect the presence of trains with sufficient precision (or other means of detecting the positions of the trains). Finally a computer system is needed that is able to reasonably predict the movements of the trains for the next few minutes.
Train speed optimization in practice:
The expensive and complicated installations usually only make sense for heavily used routes.
Swiss Federal Railways: Lenzburg-Killwangen (since 2000) Zurich-Altstetten Around Olten (installed 2004, current status unknown) Probably other places (Zürich S-Bahn?) | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Fab Tree Hab**
Fab Tree Hab:
The Fab Tree Hab is a hypothetical ecological home design developed at MIT in the early 2000s by Mitchell Joachim, Javier Arbona and Lara Greden. With the idea of easing the burden humanity places on the environment with conventional housing by growing "living, breathing" tree homes.It would be built by allowing native trees to grow over a computer-designed (CNC) removable plywood scaffold. Once the plants are interconnected and stable, the plywood would be removed and reused. MIT is experimenting with trees that grow quickly and develop an interwoven root structure that's soft enough to "train" over the scaffold, but then hardens into a more durable structure. The inside walls would be conventional clay and plaster.
Fab Tree Hab:
An old methodology new to buildings is introduced in this design: pleaching. Pleaching is a method of weaving together tree branches to form living archways, lattices, or screens. The technique is also named "aeroponic culture". The load-bearing part of the structure uses trees that self-graft or inosculate such as live oak, elm and dogwood. The lattice frame for the walls and roof are created with the branches of the trees. Vines create a dense protective layer woven along the exterior, interspersed with soil pockets and growing plants.This building could be very sustainable as it can use bio-waste for manure for the trees. Which can use the grey water from the home for the trees and garden. There are plans to be able to use rainwater. These building would improve the quality of life by giving back to nature instead of just exploiting it. Throughout the whole life cycle of the home it remains part of the ecology, feeding different organism at different times of its life. The expected life span is greater than standard structures of brick and concrete. The whole community and the individual would befit from this life style.The Fab Tree Hab is an experiment that would develop over time. Extra operating costs required over the life-time of the home include pest management with organic pesticides and maintenance of the living machine's water treatment system. Technical demonstration and innovation is still needed for certain components, primarily the bioplastic windows that accept growth of the structure and the management of flows across the wall section to assure that the interior remains dry and animal-free. All in all, the elapsed time to reach livability is greater than the traditional sense, but so should be the health and longevity of the home and family. Above all, building this home could be achieved at a minimal price. Depending on the surrounding climate the house is to be grown in, the team expect it will take a minimum of five years to complete its structure. Realization of these homes will begin as an experiment, and it is envisioned that thereafter, the concept of renewal will take on a new architectural form, one of inter-dependency between nature and people.
Fab Tree Hab:
As of May 2007 Mitchell Joachim stated that there is a "50 per cent" organic project in California, combining natural elements and traditional construction.
Trees:
Main trees suggested to be used are elms, dogwood and oaks. The teams hopes the homes can be grown using mainly native trees. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Neo-Freudianism**
Neo-Freudianism:
Neo-Freudianism is a psychoanalytic approach derived from the influence of Sigmund Freud but extending his theories towards typically social or cultural aspects of psychoanalysis over the biological.The neo-Freudian school of psychiatrists and psychologists were a group of loosely-linked American theorists/writers of the mid-20th century "who attempted to restate Freudian theory in sociological terms and to eliminate its connections with biology."
Dissidents and post-Freudians:
Dissidents The term neo-Freudian is sometimes loosely (but inaccurately) used to refer to those early followers of Freud who at some point accepted the basic tenets of Freud's theory of psychoanalysis but later dissented from it. "The best-known of these dissenters are Alfred Adler and Carl Jung.… The Dissidents."An interest in the social approach to psychodynamics was the major theme linking the so-called neo-Freudians: Alfred Adler had perhaps been "the first to explore and develop a comprehensive social theory of the psychodynamic self.": 61 Following "Adler's death, some of his views…came to exert considerable influence on the neo-Freudian theory." Indeed, it has been suggested of "Horney and Sullivan ... that these theorists could be more accurately described as 'neo-Adlerians' than 'neo-Freudians'.": 54 Post-Freudians The Independent Analysts Group of the British Psycho-Analytical Society ("Contemporary Freudians") are—like the ego-psychologists (e.g. Heinz Hartmann) or the intersubjective analysts in the States—perhaps best thought of as "different schools of psychoanalytic thought," or as "Post-Freudians…post-Freudian developments." They are distinct from the Kleinian schools of thought and include figures such as Christopher Bollas, D. W. Winnicott, and Adam Phillips.
Neo-Freudian ideas:
History As early as 1936, Erich Fromm had been independently regretting that psychoanalysts "did not concern themselves with the variety of life experience…and therefore did not try to explain psychic structure as determined by social structure." Karen Horney, too, "emphasised the role culture exerts in the development of personality and downplayed the classical driven features outlined by Freud.": 61 Erik H. Erikson, for his part, stressed that "psychoanalysis today is…shifting its emphasis…to the study of the ego's roots in the social organisation," and that its method should be "what H. S. Sullivan called 'participant', and systematically so."Doctor and psychotherapist Harald Schultz-Hencke (1892–1953) was thoroughly busy with questions like impulse and inhibition and with the therapy of psychoses as well as the interpretation of dreams. He worked with Matthias Göring in his institute (Deutsches Institut für psychologische Forschung und Psychotherapie), and created the name Neopsychoanalyse in 1945. The "Neo-Freudian revolt against the orthodox theory of instincts" was thus anchored in a sense of what Harry Stack Sullivan termed "our incredibly culture-ridden life." By their writings, and "in accessible prose, Fromm, Horney, and others mounted a cultural and social critique which became almost conventional wisdom."Through informal and more formal institutional links, such as the William Alanson White Institute, as well as through likeness of ideas, the neo-Freudians made up a cohesively distinctive and influential psychodynamic movement.
Neo-Freudian ideas:
Basic anxiety Karen Horney theorized that to deal with basic anxiety, the individual has three options: Moving toward others: Accept the situation and become dependent on others. This strategy may entail an exaggerated desire for approval or affection.
Moving against others: Resist the situation and become aggressive. This strategy may involve an exaggerated need for power, exploitation of others, recognition, or achievement.
Moving away from others: Withdraw from others and become isolated. This strategy may involve an exaggerated need for self-sufficiency, privacy, or independence.
Neo-Freudian ideas:
Basic personality The neo-Freudian Abram Kardiner was primarily interested in learning how a specific society acquires adaptation concerning its environment. He does this by forming within its members what he names a "basic personality." The "basic personality" can initially be traced to the operation of primary institutions. It ultimately creates clusters of unconscious motivations in the specific individual "which in turn are projected in the form of secondary institutions," such as reality systems. The basic personality finds expression in the secondary institutions.
Criticism:
"Fenichel developed a stringent theoretical critique of the neo-Freudians", which informed and fed into the way "Herbert Marcuse, in his 'Critique of Neo-Freudian Revisionism'...icily examines the tone of uplift and the Power of Positive Thinking that pervades the revisionists' writings, and mocks their claims to scientific seriousness."In comparable fashion, "an article…by Mr Edward Glover, entitled Freudian or Neo-Freudian, directed entirely against the constructions of Mr Alexander" equally used the term as a form of orthodox reproach.
Criticism:
In the wake of such contemporary criticism, a "consistent critique levelled at most theorists cited above is that they compromise the intrapersonal interiority of the psyche;" but one may accept nonetheless that "they have contributed an enduring and vital collection of standpoints relating to the human subject.": 66
Influence, successors, and offshoots:
In 1940, Carl Rogers had launched what would become person-centred psychotherapy, "crediting its roots in the therapy of Rank...& in the neo-Freudian analysts—especially Karen Horney.": 109 A decade later, he would report that it had "developed along somewhat different paths than the psychotherapeutic views of Horney or Sullivan, or Alexander and French, yet there are many threads of interconnection with these modern formulations of psychoanalytic thinking.": 279 A half-century further on, whether by direct or by indirect influence, "consistent with the traditions of these schools, current theorists of the social and psychodynamic self are working in the spaces between social and political theory and psychoanalysis" once again.
Influence, successors, and offshoots:
Cultural offshoots In his skit on Freud's remark that "if my name were Oberhuber, my innovations would have found far less resistance," Peter Gay, considering the notional eclipse of "Oberhuber" by his replacement Freud, adjudged that "the prospect that deviants would have to be called neo-Oberhuberians, or Oberhuberian revisionists, contributed to the master's decline."
Neo-Freudians:
Alfred Adler Erik Erikson Erich Fromm Frieda Fromm-Reichmann Karen Horney Carl Jung Abram Kardiner Harald Schultz-Hencke Harry Stack Sullivan Clara Thompson Others with possible neo-Freudian links Franz Alexander Jessica Benjamin Nancy Chodorow Richard Hakim Thomas Ogden David Rapaport Alex Unger | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**French curve**
French curve:
A French curve is a template usually made from metal, wood or plastic composed of many different segments of the Euler spiral (aka the clothoid curve). It is used in manual drafting and in fashion design to draw smooth curves of varying radii. The curve is placed on the drawing material, and a pencil, knife or other implement is traced around its curves to produce the desired result. They were invented by the German mathematician Ludwig Burmester and are also known as Burmester (curve) set.
Clothing design:
French curves are used in fashion design and sewing alongside hip curves, straight edges and right-angle rulers. Commercial clothing patterns can be personalized for fit by using French curves to draw neckline, sleeve, bust and waist variations. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Membraneless Fuel Cells**
Membraneless Fuel Cells:
Membraneless Fuel Cells convert stored chemical energy into electrical energy without the use of a conducting membrane as with other types of Fuel Cells. In Laminar Flow Fuel Cells (LFFC) this is achieved by exploiting the phenomenon of non-mixing laminar flows where the interface between the two flows works as a proton/ion conductor. The interface allows for high diffusivity and eliminates the need for costly membranes. The operating principles of these cells mean that they can only be built to millimeter-scale sizes. The lack of a membrane means they are cheaper but the size limits their use to portable applications which require small amounts of power.
Membraneless Fuel Cells:
Another type of membraneless fuel cell is a Mixed Reactant Fuel Cell (MRFC). Unlike LFFCs, MRFCs use a mixed fuel and electrolyte, and are thus not subject to the same limitations. Without a membrane, MRFCs depend on the characteristics of the electrodes to separate the oxidation and reduction reactions. By eliminating the membrane and delivering the reactants as a mixture, MRFCs can potentially be simpler and less costly than conventional fuel cell systems.The efficiency of these cells is generally much higher than modern electricity producing sources. For example, a fossil fuel power plant system can achieve a 40% electrical conversion efficiency while an outdated nuclear power plant is slightly lower at 32%. GenIII and GenIV Nuclear Fission plants can get up to 90% efficient if using direct conversion or up to 65% efficient if using a magnetohydrodynamic generator as a topping cycle{{Citation needed|reason=again, the numbers seem way off. The best achieved efficiency for initial cycle is about 30%. The capture of residual thermal energy is at best 30% to date, which comes to overall efficiency of 51% at best |date=June 2022}}. Fuel cell systems are capable of reaching efficiencies in the range of 55%–70%. However, as with any process, fuel cells also experience inherent losses due to their design and manufacturing processes.
Overview:
A fuel cell consists of an electrolyte which is placed in between two electrodes – the cathode and the anode. In the simplest case, hydrogen gas passes over the cathode, where it is decomposed into hydrogen protons and electrons. The protons pass through the electrolyte (often NAFION – manufactured by DuPont) across to the anode to the oxygen. Meanwhile, the free electrons travel around the cell to power a given load and then combine with the oxygen and hydrogen at the anode to form water. Two common types of electrolytes are a proton exchange membrane(PEM) (also known as Polymer Electrolyte Membrane) and a ceramic or solid oxide electrolyte (often used in Solid oxide fuel cells). Although hydrogen and oxygen are very common reactants, a plethora of other reactants exist and have been proven effective.
Overview:
Hydrogen for fuel cells can be produced in many ways. The most common method in the United States (95% of production) is via Gas reforming, specifically using methane, which produces hydrogen from fossil fuels by running them through a high temperature steam process. Since fossil fuels are primarily composed of carbon and hydrogen molecules of various sizes, various fossil fuels can be utilized. For example, methanol, ethanol, and methane can all be used in the reforming process. Electrolysis and high temperature combination cycles are also used to provide hydrogen from water whereby the heat and electricity provide sufficient energy to disassociate the hydrogen and oxygen atoms.
Overview:
However, since these methods of hydrogen production are often energy and space intensive, it is often more convenient to use the chemicals directly in the fuel cell. Direct Methanol Fuel Cells (DMFC's), for example, use methanol as the reactant instead of first using reformation to produce hydrogen. Although DMFC's are not very efficient (~25%), they are energy dense which means that they are quite suitable for portable power applications. Another advantage over gaseous fuels, as in the H2-O2 cells, is that liquids are much easier to handle, transport, pump and often have higher specific energies allowing for greater power extraction. Generally gases need to be stored in high pressure containers or cryogenic liquid containers which is a significant disadvantage to liquid transport.
Membraneless Fuel Cells and Operating Principles:
The majority of fuel cell technologies currently employed are either PEM or SOFC cells. However, the electrolyte is often costly and not always completely effective. Although hydrogen technology has significantly evolved, other fossil fuel based cells (such as DMFC's) are still plagued by the shortcomings of proton exchange membranes. For example, fuel crossover means that low concentrations need to be used which limits the available power of the cell. In solid oxide fuel cells, high temperatures are needed which require energy and can also lead to quicker degradation of materials. Membraneless fuel cells offer a solution to these problems.
Membraneless Fuel Cells and Operating Principles:
Laminar Flow LFFC's overcome the problem of unwanted crossover through the manipulation of the Reynolds number, which describes the behavior of a fluid. In general, at low Reynolds numbers, flow is laminar whereas turbulence occurs at a higher Reynolds number. In laminar flow, two fluids will interact primarily through diffusion which means mixing is limited. By choosing the correct fuel and oxidizing agents in LFFC's, protons can be allowed to diffuse from the anode to the cathode across the interface of the two streams. The LFFC's are not limited to a liquid feed and in certain cases, depending on the geometry and reactants, gases can also be advantageous. Current designs inject the fuel and oxidizing agent into two separate streams which flow side by side. The interface between the fluids acts as the electrolytic membrane across which protons diffuse. Membraneless fuel cells offer a cost advantage due to the lack of the electrolytic membrane. Further, a decrease in crossover also increases fuel efficiency resulting in higher power output.
Membraneless Fuel Cells and Operating Principles:
Diffusion Diffusion across the interface is extremely important and can severely affect fuel cell performance. The protons need to be able to diffuse across both the fuel and the oxidizing agent. The diffusion coefficient, a term which describes the ease of diffusion of an element in another medium, can be combined with Fick's laws of diffusion which addresses the effects of a concentration gradient and distance over which diffusion occurs: J=−D∂ϕ∂x where J is the diffusion flux in dimensions of [(amount of substance) length−2 time−1], example (molm2⋅s) . J measures the amount of substance that will flow through a small area during a small time interval.
Membraneless Fuel Cells and Operating Principles:
D is the diffusion coefficient or diffusivity in dimensions of [length2 time−1], example (m2s) ϕ (for ideal mixtures) is the concentration in dimensions of [(amount of substance) length−3], example (molm3) x is the diffusion length i.e. the distance over which diffusion occursIn order to increase the diffusion flux, the diffusivity and/or concentration need to be increased while the length needs to be decreased. In DMFC's for example, the thickness of the membrane determines the diffusion length while the concentration is often limited due to crossover. Thus, the diffusion flux is limited. A membraneless fuel cell is theoretically the better option since the diffusion interface across both fluids is extremely thin and using higher concentrations does not result in a drastic effect on crossover.
Membraneless Fuel Cells and Operating Principles:
In most fuel cell configurations with liquid feeds, the fuel and oxidizing solutions almost always contain water which acts as a diffusion medium. In many hydrogen-oxygen fuel cells, the diffusion of oxygen at the cathode is rate limiting since the diffusivity of oxygen in water is much lower than that of hydrogen. As a result, LFFC performance can also be improved by not using aqueous oxygen carriers.
Membraneless Fuel Cells and Operating Principles:
Research and development The promise of membraneless fuel cells has been offset by several problems inherent to their designs. Ancillary structures are one of the largest obstacles. For example, pumps are required to maintain laminar flow while gas separators can be needed to supply the correct fuels into the cells. For micro fuel cells, these pumps and separators need to be miniaturized and packaged into a small volume (under 1 cm3). Associated with this process is a so-called "packaging penalty" which results in higher costs. Further, pumping power drastically increases with decreasing size (see Scaling Laws) which is disadvantageous. Efficient packaging methods and/or self-pumping cells (see Research and Development) need to be developed to make this technology viable. Also, while using high concentrations of specific fuels, such as methanol, crossover still occurs. This problem can be partially solved by using a nanoporous separator, lowering fuel concentration or choosing reactants which have a lower tendency towards crossover.
Membraneless Fuel Cells and Operating Principles:
Date: January 2010: Researchers developed a novel method of inducing self-pumping in a membraneless fuel cell. Using formic acid as a fuel and sulfuric acid as an oxidant, CO2 is produced in the reaction in the form of bubbles. The bubbles nucleate and coalesce on the anode. A check valve at the supply end prevents any fuel entering while the bubbles are growing. The check valve is not mechanical but hydrophobic in nature. By creating micro structures which form specific contact angles with water, fuel cannot be drawn backwards. As the reaction continues, more CO2 is formed while fuel is consumed. The bubble begins to propagate towards the outlet of the cell. However, before the outlet, a hydrophobic vent allows the carbon dioxide to escape while simultaneously ensuring other byproducts (such as water) do not clog the vent. As the carbon dioxide is being vented, fresh fuel is also drawn in at the same through the check valve and the cycle begins again. Thus, the fuel cell pumping is regulated by the reaction rate. This type of cell is not a two stream laminar flow fuel cell. Since the formation of bubbles can disrupt two separate laminar flows, a combined stream of fuel and oxidant was used. In laminar conditions, mixing will still not occur. It was found that using selective catalysts (i.e. Not platinum) or extremely low flow rates can prevent crossover.
Scaling Issues:
Membraneless fuel cells are currently being manufactured on the micro scale using fabrication processes found in the MEMS/NEMS area. These cell sizes are suited for the small scale due to the limit of their operating principles. The scale-up of these cells to the 2–10 Watt range has proven difficult since, at large scales, the cells cannot maintain the correct operating conditions.
Scaling Issues:
For example, laminar flow is a necessary condition for these cells. Without laminar flow, crossover would occur and a physical electrolytic membrane would be needed. Maintaining laminar flow is achievable on the macro scale but maintaining a steady Reynolds number is difficult due to variations in pumping. This variation causes fluctuations at the reactant interfaces which can disrupt laminar flow and affect diffusion and crossover. However, self-pumping mechanisms can be difficult and expensive to produce on the macro-scale. In order to take advantage of hydrophobic effects, the surfaces need to be smooth to control the contact angle of water. To produce these surfaces on a large scale, the cost will significantly increase due to the close tolerances which are needed. Also, it is not evident whether using a carbon-dioxide based pumping system on the large scale is viable.
Scaling Issues:
Membraneless fuel cells can utilize self-pumping mechanisms but requires the use of fuel which release GHG's (greenhouse gases) and other unwanted products. To use an environmentally friendly fuel configuration (such as H2-O2), self pumping can be difficult. Thus, external pumps are required. However, for a rectangular channel, the pressure required increases proportional to the L−3, where L is a length unit of the cell. Thus, by decreasing the size of a cell from 10 cm to 1 cm, the required pressure will increase by 1000. For micro fuel cells, this pumping requirement requires high voltages. Although in some cases, Electroosmotic flow can be induced. However, for liquid mediums, high voltages are also required. Further, with decreasing size, surface tension effects also become significantly more important. For the fuel cell configuration with a carbon dioxide generating mechanism, the surface tension effects could also increase the pumping requirements drastically.
Potential Applications of LFFCs:
The thermodynamic potential of a fuel cell limits the amount of power that an individual cell can deliver. Therefore, in order to obtain more power, fuel cells must be connected in series or parallel (depending on whether greater current or voltage is desired). For large scale building and automobile power applications, macro fuel cells can be used because space is not necessarily the limiting constraint. However, for portable devices such as cell phones and laptops, macro fuel cells are often inefficient due to their space requirements lower run times. LFFCs however, are perfectly suited for these types of applications. The lack of a physical electrolytic membrane and energy dense fuels that can be used means that LFFC's can be produced at lower costs and smaller sizes. In most portable applications, energy density is more important than efficiency due to the low power requirements. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Underwater survey**
Underwater survey:
An underwater survey is a survey performed in an underwater environment or conducted remotely on an underwater object or region. Survey can have several meanings. The word originates in Medieval Latin with meanings of looking over and detailed study of a subject. One meaning is the accurate measurement of a geographical region, usually with the intention of plotting the positions of features as a scale map of the region. This meaning is often used in scientific contexts, and also in civil engineering and mineral extraction. Another meaning, often used in a civil, structural, or marine engineering context, is the inspection of a structure or vessel to compare actual condition with the specified nominal condition, usually with the purpose of reporting on the actual condition and compliance with, or deviations from, the nominal condition, for quality control, damage assessment, valuation, insurance, maintenance, and similar purposes. In other contexts it can mean inspection of a region to establish presence and distribution of specified content, such as living organisms, either to establish a baseline, or to compare with a baseline.
Underwater survey:
These types of survey may be done in or of the underwater environment, in which case they may be referred to as underwater surveys, which may include bathymetric, hydrographic, and geological surveys, archaeological surveys, ecological surveys, and structural or vessel safety surveys. In some cases they can be done by remote sensing, using a variety of tools, and sometimes by direct human intervention, usually by a professional diver. Underwater surveys are an essential part of the planning, and often of quality control and monitoring, of underwater construction, dredging, mineral extraction, ecological monitoring, and archaeological investigations. They are often required as part of an ecological impact study.
Types:
The types of underwater survey include, but are not necessarily restricted to, archeological, bathymetric and hydrographic, ecological, geological, and construction site surveys, and inspection surveys of marine and coastal structures and vessels afloat. A survey of the vessel structural condition and the adjacent site and hydrographic conditions would also be done when assessing proposed marine salvage operations.
Types:
Archaeological surveys Archaeological surveys of underwater sites have traditionally been done by divers, but at sites where the depth is too great, sonar surveys have been done from surface and submersible vehicles, and photomosaic techniques have been done using ROUVs. Traditional methods include direct measurement from a baseline or grid set up at the site, and triangulation by direct measurement from marks of known position installed at the site, in the same way these would be used at a terrestrial site. Accuracy may be compromised by water conditions.
Types:
This work is usually done by archaeologists who are qualified scientific divers.
Types:
Bathymetric and hydrographic surveys Bathymetric surveys are traditionally done from the surface, by measuring depth (soundings) at measured positions along transect lines and later plotting the data onto a bathymetric chart, on which lines of constant depth (isobaths) may be drawn by interpolation of soundings. It is also conventional to provide a representative set of spot depths on the chart. Originally, soundings were made manually by measuring the length of a weighted line lowered to the bottom, bur after the development of accurate and reliable echo-sounding equipment it became the standard method. Data recording was automated when the equipment became available, and later precise position data was integrated into the data sets. Multibeam sonar with GPS position data corrected for vessel motion and combined in real time is the state of the art in the early 21st century. Bathymetric surveys of some bodies of water have required different procedures, particularly for sinkholes, caverns and caves where a significant portion of the bottom walls, and in some cases ceilings, are not visible to the sounding equipment from the surface, and it has been necessary to use remotely operated underwater vehicles or divers to gather the data. One of the complications of this class of underwater survey is the relative difficulty of establishing a baseline, or an accurate position for the ROUV, as GPS signals do not propagate through water. In some cases a physical line has been used, but sometimes a baseline can be established using sonar transducers set up at accurately surveyed positions, and relative offsets measured.
Types:
Ecological surveys Various techniques have been used for underwater ecological surveys. Divers are frequently used to collect data, either by direct observation and recording, or by photographic recording at recorded locations, which may be specified to a given precision depending on the requirements of the project and available location technology.
Types:
One method is for divers to use geolocated photographs taken by divers following a route recorded by a towed surface GPS receiver on a float kept above the camera by line tension. Date and time data are recorded concurrently by the camera and GPS unit, allowing position data for each photo to be extracted by post-processing or inspection. GPS precision may be augmented by Wide Area Augmentation System (WAAS). Depth data may be captured on camera from dive computers or depth gauges carried by the divers or mounted in view of the camera. The photos may be viewed on a map or via a geographic information system (GIS) for analysis. This method can also be used for spatial surveys of small areas, particularly in places where a survey vessel cannot go. To map an area the diver tows the float along bottom contours and the GPS track is used to create a map using drafting or GIS software. Spot depths may also be taken, using a digital camera to record time and depth from a depth gauge or dive computer to synchronize with the track data. This procedure can be combined with photographic recording of the benthic communities at intervals along the contour or perimeter.
Types:
Surveys by professional divers tend to be relatively expensive, and some ecological monitoring programs and data gathering programs have enlisted the aid of volunteer recreational divers to conduct data collection appropriate to their certification and in some cases, further training, such as the Australian-based Reef Life Survey. Others, such as iNaturalist, have used the crowdsourcing system of uploaded digital photographic records of observations, with location data to whatever standard is available, which can vary considerably, thereby taking advantage of the thousands of amateur photographers who record their underwater surroundings anyway. In this way millions of observations from dive sites all over the world have been accumulated.Types of ecological survey: Transect – Path along which the observer counts and records occurrences of the subjects of the survey Quadrat – Rectangular frame used to demarcate a part of the substrate for detailed analysis Photogrammetry – Taking measurements using photography Baited remote underwater video – Equipment for estimating fish populationsSometimes more than one type of observations are combined in a survey. For example, the Reef Life Survey procedure includes three components along the same transect: Visual count of fish, visual count of benthic fauna, and photographs of the bottom at regular intervals.
Types:
Geological surveys A geological survey is the systematic investigation of the geology beneath a given piece of ground for the purpose of creating a geological map or model. Underwater geological surveying employs techniques from the underwater equivalent of a traditional walk-over survey, studying outcrops and landforms, to intrusive methods, such as boreholes, to the use of geophysical techniques and remote sensing methods. An underwater geological survey map typically superimposes the surveyed extent and boundaries of geological units on a bathymetric map, together with information at points (such as measurements of orientation of bedding planes) and lines (such as the intersection of faults with the seabed surface). The map may include cross sections to illustrate the three-dimensional interpretation. Much of this work is done from surface vessels by remote sensing, bur in some cases such as in flooded caves, measurement and sampling requires remotely operated underwater vehicles or direct intervention by divers.
Types:
Reflection seismology techniques are used for shipborne subsurface remote sensing. Seismic sources include air guns, sparkers and boomers.
Airborne geophysical methods include magnetic, electromagnetic, and gravity measurement.
Types:
Site surveys Site surveys are inspections of an area where work is proposed, to gather information for a design. It can determine a precise location, access, best orientation for the site and the location of obstacles. The type of site survey and the best practices required depend on the nature of the project. In hydrocarbon exploration, for example, site surveys are run over the proposed locations of offshore exploration or appraisal wells. They consist typically of a tight grid of high resolution (high frequency) reflection seismology profiles to look for possible gas hazards in the shallow section beneath the seabed and detailed bathymetric data to look for possible obstacles on the seafloor (e.g. shipwrecks, existing pipelines) using multibeam echosounders.
Types:
A type of site survey is performed during marine salvage operations, to assess the structural condition of a stranded vessel and to identify aspects of the vessel, site and environment that may affect the operation. Such a survey may include investigation of hull structural and watertight integrity, extent of flooding, bathymetry and geology of the immediate vicinity, currents and tidal effects, hazards, and possible environmental impact of the salvage work.
Types:
Structural surveys Structural integrity inspections of inland, coastal and offshore underwater structures, including bridges, dams, causeways, harbours, breakwaters, jetties, embankments, levees, petroleum and gas production platforms and infrastructure, pipelines, wellheads and moorings.
Types:
Vessel safety surveys Vessel safety surveys are inspections of the structure and equipment of a vessel to assess the condition of the surveyed items and check that they comply with legal or classification society requirements for insurance and registration. They may occur at any time when there is reason to suspect that the condition has changed significantly since the previous survey, or as a condition of purchase, and the first survey is generally during construction (built under survey) or before first registration. The criteria for acceptance are defined by the licensing or registration authority for a variety of equipment vital to the safe operation of the vessel, such as hull structure, static stability, propulsion machinery, auxiliary machinery, safety equipment, lifting equipment, rigging, ground tackle, etc.
Types:
Some surveys must be done in dry dock, but this is expensive, and in some cases for intermediate surveys the underwater part of the external survey may be done afloat using divers or ROUVs to do the inspection, usually providing live video to the surveyor, or possibly video recording for later analysis. Live video has the advantage that the surveyor can instruct the diver to investigate further or provide views from other angles. Live video would normally also be recorded for the records.
Tools:
Remote measurement through water Single beam echosounders are used to measure distance of a reflecting surface, like the seabed, by comparing the time between emission of a sound signal and first receiving the reflected signal back at the transceiver, using the speed of sound in water. They are usually used to make a series of spot depth measurements along the path of the transducer, which can be used to map the bottom profile.
Tools:
Multibeam echosounders use beamforming to extract directional information from the returning sound waves, producing a swath of depth readings across the path of the transducer from a single ping. The rate of data acquisition is far greater than for single beam systems, but they are susceptible to shadowing effects from high-profile surfaces offset to the side of the transducer path. This can be compensated by overlapping swaths. The data is processed to give a three dimensional image of the bottom.
Tools:
Acoustic Doppler current profilers (ADCP) are hydro-acoustic current meters, used to measure water current velocities over a depth range using the Doppler effect of sound waves scattered back from particles within the water column. The traveling time of the sound waves gives the distance, and the frequency shift of the echo is proportional to the water velocity along the acoustic path.
Tools:
Lidar uses a laser light source and optical receiver to measure range and direction of reflected signals, but is limited by the water transparency.
Side-scan sonar is used to efficiently create images of large areas of the sea floor, as seen from the point of view of the transducer.
Tools:
Seismic sources such as sparkers and boomers are used in seismic reflection profiling, using sound pulse frequencies that effectively penetrate the solid seabed and are partially reflected by changes in acoustic impedance, often signifying a change in rock type. Boomers work in the 500 to 4000 Hz range. and sparkers in the 200 to 800 Hz range. Lower frequency will usually penetrate to greater depth, but with lower resolution.
Tools:
Platforms Dedicated survey vessels and vessels of opportunity. Diving support vessels for surface-supplied diving operations, and dive boats for scuba surveys. DSVs are often fitted for ROV support and other underwater surveys.
Autonomous survey vessels are more economical to operate than manned vessels, and can be sent into waters that are too shallow or confined or otherwise hazardous for larger manned vessels.
Tools:
Autonomous underwater vehicles are more economical than manned vehicles. Researchers have focused on the development of AUVs for long-term data collection in oceanography and coastal management. The oil and gas industry uses AUVs to make detailed maps of the seafloor before they start building subsea infrastructure. The AUV allows survey companies to conduct precise surveys of areas where traditional bathymetric surveys would be less effective or too costly. Also, post-lay pipe surveys which include pipeline inspection are possible. The use of AUVs for pipeline inspection and inspection of underwater man-made structures is becoming more common. Scientists use AUVs to study lakes, the ocean, and the ocean floor. A variety of sensors can also be carried to measure the concentration of various elements or compounds in the water, the absorption or reflection of light, and the presence of microscopic life. Examples include conductivity-temperature-depth sensors (CTDs), chlorophyll fluorometers, and pH sensors.
Tools:
Remotely operated underwater vehicles. Survey or inspection ROVs are generally smaller than workclass ROVs and are often sub-classified as either Class I: Observation Only or Class II Observation with payload. They are used to assist with hydrographic survey, and also for inspection work. Survey ROVs, although smaller than workclass, often have comparable performance with regard to the ability to hold position in currents, and often carry similar tools and equipment - lighting, cameras, sonar, USBL (Ultra-short baseline) beacon, and strobe flasher depending on the payload capability of the vehicle and the needs of the user.
Tools:
Underwater position measurement systems Underwater acoustic positioning systems are systems for the tracking, navigation and location of underwater vehicles or divers by means of acoustic distance and/or direction measurements, and subsequent position triangulation. They are commonly used in a wide variety of underwater work, including oil and gas exploration, ocean sciences, salvage operations, marine archaeology, law enforcement and military activities.
Tools:
Long baseline acoustic positioning systems (LBL systems) use networks of sea-floor mounted baseline transponders as reference points for navigation. These are generally deployed around the perimeter of a work site. The LBL technique results in very high positioning accuracy and position stability that is independent of water depth. It is generally better than 1-meter and can reach a few centimeters accuracy. LBL systems are generally used for precision underwater survey work where the accuracy or position stability of ship-based short or ultra-short baseline positioning systems does not suffice.
Tools:
Short baseline acoustic positioning system (SBL acoustic positioning systems) SBL systems do not require any seafloor mounted transponders or equipment and are thus suitable for tracking underwater targets from boats or ships that are either anchored or under way. However, unlike USBL systems, which offer a fixed accuracy, SBL positioning accuracy improves with transducer spacing. Thus, where space permits, such as when operating from larger vessels or a dock, the SBL system can achieve a precision and position robustness that is similar to that of sea floor mounted LBL systems, making the system suitable for high-accuracy survey work. When operating from a smaller vessel where transducer spacing is limited (i.e. when the baseline is short), the SBL system will exhibit reduced precision.
Tools:
Ultra-short baseline acoustic positioning system (USBL), also known as super short base line (SSBL), consists of a transceiver, which is mounted under a ship, and a transponder or responder on the seafloor, on a towfish, or on an ROV. A computer, is used to calculate a position from the ranges and bearings measured by the transceiver. USBLs are also used in "inverted" (iUSBL) configurations, with the transceiver mounted on an autonomous underwater vehicle, and the transponder on the installation that launches it. In this case, the signal processing happens inside the vehicle to allow it to locate the transponder for applications such as automatic docking and target tracking.
Tools:
Manual measurement underwater Depth measurement: Depth gauge (using pressure as a proxy) Dive computer (using pressure as a proxy) Measuring tape (direct linear measurement) Pneumofathometer (using pressure as a proxy)Length measurement in other directions: Measuring tape Surveyor's chain Calibrated distance line Towed GPS receiver on float (by spherical trigonometry) Derived by triangulation from a baseline, angular measurement, and trigonometry Inertial navigation system, integrated from accelerometer output in three dimensions Hand-held range-finding sonar Vernier and plain calipersAngular measurement: Magnetic compass Protractor Clinometer Goniometer Derived from GPS positions Derived from linear triangulation and trigonometry Derived from inertial navigation position dataNon-destructive testing: Ultrasonic thickness measurement Sampling and specimen collection Samples of seafloor sediments and rock can be collected using grabs, coring devices, ROUVs and divers. Coring devices include core drills and impact penetrators. Divers and ROUV operators are more discriminating in their selection of samples than grabs and remotely operated coring devices. Biological samples can be collected by dredges, grabs, traps, or nets, but more directed sampling generally requires visual input and human intervention, and is commonly done by divers, ROUVs and manned submersibles equipped for collection.
Tools:
Recording and counting Underwater photography. Digital underwater cameras can conveniently be used to record an image, and the time at which the photo was taken. In some cases direction, inclination and depth are also available from the camera, or can be recorded by photographing the display of appropriate instruments.
Jump cameras are cameras mounted on a frame that triggers an exposure when the frame hits the bottom. To operate, the frame is lowered until the rope slacks off, then lifted and the boat moved to the next position.
Underwater videography, the branch of electronic underwater photography concerned with capturing underwater moving images, and live video feeds, which allow a remote operator to see the underwater environment from elsewhere.
Tools:
Baited remote underwater video (BRUV) is a system used in marine biology research. By attracting fish into the field of view of a remotely controlled camera, the technique records fish diversity, abundance and behaviour of species. Sites are sampled by video recording the region surrounding a baited canister which is lowered to the bottom. The video can be transmitted directly to the surface by cable, or recorded for later analysis. Baited cameras are highly effective at attracting scavengers and subsequent predators, and are a non-invasive method of generating relative abundance indices for a number of marine species.
Tools:
Checklists are useful when a reasonably small range of objects or types of object are to be recorded as being present, as they reduce the amount of writing that must be done underwater. (legibility tends to suffer in cold water or in moving water) Clipboard or diver's slate and pencil, are used when sketches and measurements are to be recorded, and are versatile though not very efficient for data recording. Waterproof paper is available for use on clipboards, and can be printed with checklists.
Tools:
Quadrat frames are used to establish a discrete area for examination, and can be visually examined, photographed, or both.
Presentation of results:
Results of underwater surveys can be presented in several ways, depending on the target demographic and intended use of the data.
A common presentation format is a map indicating spatial distribution or general topography, often involving a depth dimension. Drawings, photographic images, graphs, tables, and text descriptions may also be used, often in conjunction with one or more maps. Maps may also be used to indicate variations over time in comparison with a baseline. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Phases of fluorine**
Phases of fluorine:
Fluorine forms diatomic molecules (F2) that are gaseous at room temperature with a density about 1.3 times that of air. Though sometimes cited as yellow-green, pure fluorine gas is actually a very pale yellow. The color can only be observed in concentrated fluorine gas when looking down the axis of long tubes, as it appears transparent when observed from the side in normal tubes or if allowed to escape into the atmosphere. The element has a "pungent" characteristic odor that is noticeable in concentrations as low as 20 ppb.
Phases of fluorine:
Fluorine condenses to a bright yellow liquid at −188 °C (−307 °F), which is near the condensation temperatures of oxygen and nitrogen. The solid state of fluorine relies on Van der Waals forces to hold molecules together, which, because of the small size of the fluorine molecules, are relatively weak. Consequently, the solid state of fluorine is more similar to that of oxygen or the noble gases than to those of the heavier halogens.
Phases of fluorine:
Fluorine solidifies at −220 °C (−363 °F) into a cubic structure, called beta-fluorine. This phase is transparent and soft, with significant disorder of the molecules; its density is 1.70 g/cm3. At −228 °C (−378 °F) fluorine undergoes a solid–solid phase transition into a monoclinic structure called alpha-fluorine. This phase is opaque and hard, with close-packed layers of molecules, and is denser at 1.97 g/cm3. The solid state phase change requires more energy than the melting point transition and can be violent, shattering samples and blowing out sample holder windows.Solid fluorine received significant study in the 1920s and 30s, but relatively less until the 1960s. The crystal structure of alpha-fluorine given, which still has some uncertainty, dates to a 1970 paper by Linus Pauling. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Linolelaidic acid**
Linolelaidic acid:
Linolelaidic acid is an omega-6 trans fatty acid (TFA) and is a cis–trans isomer of linoleic acid. It is found in partially hydrogenated vegetable oils. It is a white (or colourless) viscous liquid.
TFAs are classified as conjugated and nonconjugated, corresponding usually to the structural elements −CH=CH−CH=CH− and −CH=CH−CH2−CH=CH−, respectively. Nonconjugated TFAs are represented by elaidic acid and linolelaidic acid. Their presence is linked heart diseases. The TFA vaccenic acid, which is of animal origin, poses less of a health risk. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Code-break procedure**
Code-break procedure:
A code-break procedure is a set of rules which determine when planned unblinding should occur in a blinded experiment. FDA guidelines recommend that sponsors of blinded trials include a code-break procedure in their standard operating procedure. A code-break procedure should only allow a participant to be unblinded before the conclusion of a trial in the event of an emergency. Code-break usually refers to the unmasking of treatment allocation, but can refer to any form of unblinding.
Code-break procedure:
Traditionally, each patient's treatment allocation data was stored in a sealed envelopes, which was to be opened to break code. However, this system is prone to abuse. Reports of researchers opening envelopes prematurely or holding the envelopes up to lights to determine their contents has led some researchers to say that the use of sealed envelopes is no longer acceptable. As of 2016, sealed envelopes were still in use in some clinical trials. Modern clinical trials usually store this information in computer files. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**ISO 22300**
ISO 22300:
ISO 22300:2021, Security and resilience – Vocabulary, is an international standard developed by ISO/TC 292 Security and resilience. This document defines terms used in security and resilience standards and includes 360 terms and definitions. This edition was published in the beginning of 2021 and replaces the second edition from 2018.
Scope and contents:
ISO 22300:2018 contains definitions for the following terms:
Related standards:
ISO 22301 Security and resilience – Business continuity management systems – Requirements | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Directory Opus**
Directory Opus:
Directory Opus (or "DOpus" as its users tend to call it) is a file manager program, originally written for the Amiga computer system in the early to mid-1990s. Commercial development on the version for the Amiga ceased in 1997. Directory Opus is still being actively developed and sold for the Microsoft Windows operating system by GPSoftware and there are open source releases of Directory Opus 4 and 5 for Amiga.
Directory Opus:
Directory Opus was originally developed by, and is still written by, Australian Jonathan Potter. Until 1994, it was published by well-known Amiga software company Inovatronics, when Potter joined with Greg Perry and the Australian-based GPSoftware to continue its development, and has since been published by GPSoftware.
Features:
Directory Opus has evolved since its first release in 1990 as a basic two-panel file manager. The interface has evolved significantly due to feedback given by its users. Some of the features include: Single or dual-panel exploring.
Folder tree (either shared or separate for dual-display).
Tabbed explorer panels.
Ability to maintain date created/modified timestamps for both files and folders.
Internal handling of ZIP, RAR, 7Zip and other archive formats (browse them like folders).
Internal FTP handling, including (for a small extra fee) advanced FTP and SSH (browse these like folders also).
Internal MTP handling for portable devices like phones and cameras.
Flat-file display, where you can flatten a folder tree and even hide the folders themselves.
Powerful file selection and renaming tools, with advanced regex.
User-definable toolbars, menus, filetypes and filetype groups.
Preview panel, with preview of thumbnails (including animated avi thumbnails).
File collections. These are like virtual folders that contain links to the original files (unlike shortcuts, these actually deal with the files directly).
History:
Release history Amiga release history Opus 1: January 1990 Opus 2: February 1991 Opus 3: 1991-12-01 Opus 4: 1992-12-04 Opus 5: 1995-04-12 Opus 5.5: 1996-08-01 Opus Magellan (5.6): 1997-05-17 Opus Magellan II (5.8): 1998-11-01 Opus Magellan II GPL (5.90): 2014-05-11Versions 1 and 2 were only available direct from the author. Versions 3 and 4 were published by Inovatronics. Versions since 5 have been published by GPSoftware (German versions were published by Stefan Ossowskis Schatztruhe). The full version of Magellan II is included for free with AmiKit package.
History:
Windows major release history Opus 6: 2001-06-18 Opus 8: 2004-10-04 Opus 9: 2007-04-27 Opus 10: 2011-04-30 Opus 11: 2014-03-03 Opus 12: 2016-09-05All Windows versions published by GPSoftware. (German versions published by Haage & Partner Computer GmbH.) Open source release history GPSoftware released the older Amiga Directory Opus 4 source code in 2000 as open-source under the GNU General Public License. AmigaOS4, AROS and MorphOS ports of this version were made available. Magellan II was released as open source under the AROS Public License in December 2012.The open source Worker (file manager) is heavily inspired by the Directory Opus 4 series. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Distal 18q-**
Distal 18q-:
Distal 18q- is a genetic condition caused by a deletion of genetic material within one of the two copies of chromosome 18. The deletion involves the distal section of 18q and typically extends to the tip of the long arm of chromosome 18.
Presentation:
Distal 18q- causes a wide range of medical and developmental concerns, with significant variation in severity due to the variation in breakpoints reported in individuals with distal 18q-. Current research is focused on establishing genotype-phenotype correlations to enable predictive genotyping.
Presentation:
Congenital anomalies Heart abnormalities are present in 25–35% of people with distal 18q-. The majority of these defects are septal. Congenital orthopedic anomalies are also relatively common, particularly rocker-bottom feet or clubfoot. Cleft lip and palate are relatively common in people with distal 18q-. Kidney abnormalities have also been reported and include horseshoe kidney, hydronephrosis, polycystic kidney, and absent kidney. Boys with distal 18q- may have genital anomalies, the most frequent being cryptorchidism and hypospadias.
Presentation:
Neurologic Hypotonia is a common finding. Around 10% of people with distal 18q- have seizures.
MRI abnormalities Dysmyelination is a common finding in people with distal 18q-, present in about 95%. Hypoplasia of the corpus callosum is also a common finding.
Vision Strabismus and nystagmus are prevalent in distal 18q-. Changes in the optic nerve, as well as colobomas, are also fairly common. Myopia has been reported in some individuals.
Presentation:
Ear and sinus infections Due to changes in facial structures, infants, toddlers, and children with distal 18q- often have poor drainage from the middle ears, leading to a build-up of fluid. This can in turn lead to recurrent ear and sinus infections. Antibiotics are typically required to treat these infections. In addition, the diagnosis of ear infections in children with 18q- is frequently complicated by stenosis or atresia of the ear canals, a common finding in people with distal 18q-.
Presentation:
Hearing People with distal 18q- frequently have conductive and/or sensorineural hearing loss. The degree of hearing loss may vary from mild to severe.
Gastrointestinal Individuals with distal 18q- may have problems with reflux. Hernias have also been reported.
Genitourinary As mentioned above, males with distal 18q- may have cryptorchidism. Hypospadias and chordee have also been reported. Also, a variety of kidney malformations have been reported in infants with distal 18q-, as noted above. Additionally, vesicouretereral reflux has been diagnosed in several people with distal 18q-.
Orthopedics As mentioned above, distal 18q- is associated with an increased incidence of clubfoot and rocker bottom feet. Also, a significant chance of developing pes planus or pes cavus exists. People with distal 18q- frequently have overlapping toes. Scoliosis and genu varum are also known orthopedic complications in children and adults with distal 18q-.
Growth Children and adults with distal 18q- are often small for their age. Many people with distal 18q- have an abnormal response to growth hormone stimulation. Those who have been treated with growth hormone have responded well to the treatment. Microcephaly is also common in people with distal 18q-.
Thyroid Hypothyroidism has been reported in some people with distal 18q-.
Immunology Several people with distal 18q- have been diagnosed with low IgA levels, resulting in an increased incidence of infections.
Presentation:
Psychiatry An increased incidence of psychiatric conditions occurs within the distal 18q- population. In one study, nearly 60% had depressive symptoms, 60% had symptoms of an anxiety disorder, 25% had manic symptoms, and 25% had psychotic symptoms. However, this study included young patients, many of whom were too young to exhibit signs of certain psychiatric conditions. The typical age of onset for many of these conditions appears to be during the teen years. Thus, the results of this study may actually underestimate the true incidence of psychiatric conditions within this population. Outbursts, or anger issues, such as temper tantrums are also common.
Presentation:
Cognition and adaptive skills 97% of individuals possess some form of intellectual disability, ranging from moderate to severe cases.The intellectual development of individuals with distal 18q- vary quite widely. In one study of 46 individuals with distal 18q-, IQ ranged from 49 to 113, with most individuals falling in the mild to moderate range of intellectual disability. Some of those with IQ scores on the lower end of the spectrum probably actually had deletions encompassing the TCF4 gene.
Presentation:
An increased incidence of autism is seen within the distal 18q- population. In a recent study, 45 of 105 individuals evaluated fell into the "possible" or "very likely" levels of risk for autism. Adaptive skills may also be delayed in people with distal 18q-.
Dysmorphology Common facial features include midfacial hypoplasia, short and downward- or upward-slanting palpebral fissures, epicanthic folds, and low-set ears with a prominent antihelix.
Genetics:
Distal 18q- is a deletion of the long arm of chromosome 18. The majority of deletions have breakpoints between 45,405,887 and the tip of the chromosome. There are no common breakpoints, thus the size of the deletions vary widely. The largest deletion reported is 30.076 Mb, while the smallest deletion reported to cause a clinical phenotype is 3.78 Mb.
Diagnosis:
Suspicion of a chromosome abnormality is typically raised due to the presence of developmental delays or birth defects. Diagnosis of distal 18q- is usually made from a blood sample. A routine chromosome analysis, or karyotype, is usually used to make the initial diagnosis, although it may also be made by microarray analysis. Increasingly, microarray analysis is also being used to clarify breakpoints. Prenatal diagnosis is possible using amniocentesis or chorionic villus sampling.
Treatment:
At present, treatment for distal 18q- is symptomatic, meaning the focus is on treating the signs and symptoms of the conditions as they arise. To ensure early diagnosis and treatment, people with distal 18q- are suggested to undergo routine screenings for thyroid, hearing, and vision problems.
History:
Distal 18q- was first described in 1964. Originally, it was called "De Grouchy syndrome" or "De Grouchy syndrome 2". Today, the preferred nomenclature for this condition is 18q-. Since this condition was originally described, researchers have clarified the size and nature of these deletions. In general, deletions of 18q fall into one of two categories: interstitial deletions, which typically have breakpoints between 18q11.2 (18.9 Mb) to 18q21.1 (43.8 Mb), and terminal deletions, which typically have a breakpoint distal to 18q21.1 (45.4 Mb) and extend to the end of the chromosome. If possible, it is preferable to indicate the general location of the deletion with the phrases "proximal 18q-" and "distal 18q-".
Research:
Currently, research is focusing on identifying the role of the genes on 18q in causing the signs and symptoms associated with distal deletions of 18q.
Research:
TCF4 – In 2007, deletions of or point mutations in this gene were identified as the cause of Pitt-Hopkins syndrome. This is the first gene that has been definitively shown to directly cause a clinical phenotype when deleted. If a deletion includes the TCF4 gene (located at 55,222,331-55,664,787), features of Pitt-Hopkins may be present, including abnormal corpus callosum, short neck, small penis, accessory and wide-spaced nipples, broad or clubbed fingers, and sacral dimple. Those with deletions inclusive of TCF4 have a significantly more severe cognitive phenotype.TSHZ1 - Point mutations and deletions of this gene are linked with congenital aural atresia. Individuals with deletions inclusive of this gene have a 78% chance of having aural atresia.
Research:
Critical regions – Recent research has narrowed the critical regions for four features of the distal 18q- phenotype down to a small segment of distal 18q, although the precise genes responsible for those features remain to be identified.The table below shows the established critical regions for four features of distal 18q-, as well as the penetrance for each of those features. The penetrance figure represents the likelihood a person would have the feature given the critical region is deleted.
Research:
Haplolethal regions - Two regions on chromosome 18 have never been found to be deleted. They are located between the centromere and 22,826,284 bp (18q11.2) and between 43,832,732 and 45,297,446 bp (18q21.1). The genes in these regions are thought to be lethal when deleted. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Emmetropia**
Emmetropia:
Emmetropia is the state of vision in which a faraway object at infinity is in sharp focus with the ciliary muscle in a relaxed state. That condition of the normal eye is achieved when the refractive power of the cornea and eye lens and the axial length of the eye balance out, which focuses rays exactly on the retina, resulting in perfectly sharp distance vision. A human eye in a state of emmetropia requires no corrective lenses for distance; the vision scores well on a visual acuity test (such as an eye chart test).While emmetropia implies an absence of myopia, hyperopia, and other optical aberrations such as astigmatism, a less strict definition requires the spherical equivalent to be between -0.5 and +0.5 D and low enough aberrations such that 20/20 vision is achieved without correction.
Emmetropia:
For example, on a Snellen chart test, emmetropic eyes score at least "6/6"(m) or "20/20"(ft) vision, meaning that at a distance of 20 ft (the first number) they see as well as a "normal" eye at a distance of 20 ft (the second number). Eyes that have enough myopia (near-sighted), hyperopia (far-sighted, excluding latent and facultative hyperopia), or optical aberrations would score worse, e.g. 20/40 (visual acuity of 0.5). Typical emmetropic vision might be 20/15 to 20/10 (visual acuity of 1.3 to 2).Emmetropes with presbyopia might use lenses for near vision.
Overview:
Emmetropia is a state in which the eye is relaxed and focused on an object more than 6 meters or 20 feet away. The light rays coming from that object are essentially parallel, and the rays are focused on the retina without effort. If the gaze shifts to something closer, light rays from the source are too divergent to be focused without effort. In other words, the eye is automatically focused on things in the distance unless a conscious effort is made to focus elsewhere. For a wild animal or human prehistorical ancestors, that arrangement would be adaptive because it allows for alertness to predators or prey at a distance.
Overview:
Accommodation of the lens does not occur in emmetropia, and the lens is about 3.6 mm thick at the center; in accommodation, it thickens to about 4.5 mm. A relatively thin lens and relatively dilated pupil are also associated. The lens usually stiffens with age, causing less ability to focus when the eyes are not in a state of emmetropia.Corrective eye surgery such as LASIK and PRK aims to correct anemmetropic vision. This is accomplished by ensuring the curvature of the cornea, the shape of the lens and their distances from each other and the retina are in harmony. By shaping the cornea, emmetropic vision can be achieved without corrective lenses. The correction for only emmetropic vision is often the reason that patients are advised to keep wearing glasses to read as they age because of presbyopia.
Emmetropization:
The development of an eye towards emmetropia is known as emmetropization. This process is guided by visual input, and the mechanisms that coordinate this process are not fully understood. It is assumed that emmetropization occurs via an active mechanism by which defocus drives growth of the eye and that genetic factors and emmetropization both influence the growth of the eye's axis. Newborns are typically hypermetropic and then undergo a myopic shift to become emmetropic.There has been some research on causal factors involved in the development of myopia and of hyperopia. In particular, prolonged near work is correlated with the development of myopia. Furthermore, outdoor activity has been found to have a protective effect on myopia development in children.
Emmetropization:
It has long been assumed that wearing corrective spectacles might possibly perturb the process of emmetropization in young children, with this assumption being supported in particular also by animal studies. However, undercorrection of myopia in humans has been shown to increase the rate of myopic progression. However, it is not yet fully understood for which patient groups, if any, the wearing of corrective spectacles in childhood actually impedes emmetropization.In hyperopic children, yet more factors are to be considered: Hyperopia is known to be a significant risk factor for esotropia, therefore undercorrection may have the side effect of increasing this risk. There is widespread consensus that undercorrection is counterindicated for children with accommodative esotropia. It is still unclear for which hyperopic, non-strabismic children corrective spectacles may translate to a lower strabismus risk. There are indications that emmetropization is relevant for hyperopic children who have at most about 3.0 diopters, whereas children with stronger hyperopia seem to not change their refraction independently of whether the refractive error is corrected or not.A Cochrane Review of three trials seeking to determine whether spectacle correction reduced the occurrence of strabismus in children included one study which suggested that spectacle correction perturbed emmetropization in children, while a second study reported no differences.
Etymology:
"Emmetropia" is derived from Greek ἔμμετρος emmetros "well-proportioned" (from ἐν en "in" and μέτρον metron "measure") and ὤψ ōps "sight" (GEN ὠπός ōpos). Translated literally, the term indicates the condition of an eye's having in itself (i.e., without recourse to corrective lenses or other instruments) the capability to obtain an accurate measurement of an object's physical appearance. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Java collections framework**
Java collections framework:
The Java collections framework is a set of classes and interfaces that implement commonly reusable collection data structures.Although referred to as a framework, it works in a manner of a library. The collections framework provides both interfaces that define various collections and classes that implement them.
Differences from Arrays:
Collections and arrays are similar in that they both hold references to objects and they can be managed as a group. However, unlike arrays, Collections do not need to be assigned a certain capacity when instantiated. Collections can grow and shrink in size automatically when objects are added or removed. Collections cannot hold primitive data types such as int, long, or double. Instead, Collections can hold wrapper classes such as java.lang.Integer, java.lang.Long, or java.lang.Double.Collections are generic and hence invariant, but arrays are covariant. This can be considered an advantage of generic objects such as Collection when compared to arrays, because under circumstances, using the generic Collection instead of an array prevents run time exceptions by instead throwing a compile-time exception to inform the developer to fix the code. For example, if a developer declares an Object[] object, and assigns the Object[] object to the value returned by a new Long[] instance with a certain capacity, no compile-time exception will be thrown. If the developer attempts to add a String to this Long[] object, the java program will throw an ArrayStoreException. On the other hand, if the developer instead declared a new instance of a Collection<Object> as ArrayList<Long>, the Java compiler will (correctly) throw a compile-time exception to indicate that the code is written with incompatible and incorrect type, thus preventing any potential run-time exceptions.The developer can fix the code by instantianting Collection<Object> as an ArrayList<Object> object. If the code is using Java SE7 or later versions, the developer can instatiate Collection<Object> as an ArrayList<> object by using the diamond operatorCollections are generic and hence reified, but arrays are not reified.
History:
Collection implementations in pre-JDK 1.2 versions of the Java platform included few data structure classes, but did not contain a collections framework. The standard methods for grouping Java objects were via the array, the Vector, and the Hashtable classes, which unfortunately were not easy to extend, and did not implement a standard member interface.To address the need for reusable collection data structures, several independent frameworks were developed, the most used being Doug Lea's Collections package, and ObjectSpace Generic Collection Library (JGL), whose main goal was consistency with the C++ Standard Template Library (STL).The collections framework was designed and developed primarily by Joshua Bloch, and was introduced in JDK 1.2. It reused many ideas and classes from Doug Lea's Collections package, which was deprecated as a result. Sun Microsystems chose not to use the ideas of JGL, because they wanted a compact framework, and consistency with C++ was not one of their goals.Doug Lea later developed a concurrency package, comprising new Collection-related classes. An updated version of these concurrency utilities was included in JDK 5.0 as of JSR 166.
Architecture:
Almost all collections in Java are derived from the java.util.Collection interface. Collection defines the basic parts of all collections. The interface has the add(E e) and remove(E e) methods for adding to and removing from a Collection respectively. It also has the toArray() method, which converts the Collection into an array of Objects in the Collection (with return type of Object[]). Finally, the contains(E e) method checks if a specified element exists in the Collection. The Collection interface is a subinterface of java.lang.Iterable, so any Collection may be the target of a for-each statement. (The Iterable interface provides the iterator() method used by for-each statements.) All Collections have an java.util.Iterator that goes through all of the elements in the Collection.
Architecture:
Collection is generic. Any Collection can store any Object. For example, any implementation of Collection<String> contains String objects. No casting is required when using the String objects from an implementation of Collection<String>. Note that the angled brackets < > can hold a type argument that specifies which type the Collection holds.
Types of collection There are several generic types of Collection: Queues, maps, lists and sets.
Queues allow the programmer to insert items in a certain order and retrieve those items in the same order. An example is a waiting list. The base interfaces for queues are called Queue.
Dictionaries/Maps store references to objects with a lookup key to access the object's values. One example of a key is an identification card. The base interface for dictionaries/maps is called Map.
Lists are finite collections where it can store the same value multiple times.
Sets are unordered collections that can be iterated and contain each element at most once. The base interface for sets is called Set.
List interface:
Lists are implemented in the collections framework via the java.util.Listinterface. It defines a list as essentially a more flexible version of an array. Elements have a specific order, and duplicate elements are allowed. Elements can be placed in a specific position. They can also be searched for within the list.
List implementations There are several concrete classes that implement List, including AbstractList and all of its corresponding subclasses, as well as CopyOnWriteArrayList.
AbstractList class The direct subclasses of AbstractList class include AbstractSequentialList, ArrayList and Vector.
AbstractList is an example of a skeletal implementation, which leverages and combines the advantages of interfaces and abstract classes by making it easy for the developer to develop their own implementation for the given interface.
ArrayList class The java.util.ArrayList class implements the List as an array. Whenever functions specific to a List are required, the class moves the elements around within the array in order to do it.
LinkedList class The java.util.LinkedList class stores the elements in nodes that each have a pointer to the previous and next nodes in the List. The List can be traversed by following the pointers, and elements can be added or removed simply by changing the pointers around to place the node in its proper place.
Vector class The Vector class has Stack as its direct subclass. This is an example of a violation of the composition over inheritance principle in the Java platform libraries, since in computer science, a vector is generally not a stack. Composition would have been more appropriate in this scenario.
Stack class The Stack class extends class java.util.Vector with five operations that allow a Vector to be treated as a Stack.
List interface:
Stacks are created using java.util.Stack. The Stack offers methods to put a new object on the Stack (method push(E e)) and to get objects from the Stack (method pop()). A Stack returns the object according to last-in-first-out (LIFO), e.g. the object which was placed latest on the Stack is returned first. java.util.Stack is a standard implementation of a stack provided by Java.
List interface:
The Stack class represents a last-in-first-out (LIFO) stack of objects. The Stack class has five additional operations that allow a Vector to be treated as a Stack. The usual push(E e) and pop() operations are provided, as well as a method (peek()) to peek at the top item on the Stack, a method to test for whether the Stack is empty (empty()), and a method to search the Stack for an item and discover how far it is from the top (search(Object o)). When a Stack is first created, it contains no items.
List interface:
CopyOnWriteArrayList class The CopyOnWriteArrayList extends the Object class, and does not extend any other classes. CopyOnWriteArrayList allows for thread-safety without performing excessive synchronization.In some scenarios, synchronization is mandatory. For example, if a method modifies a static field, and the method must be called by multiple threads, then synchronization is mandatory and concurrency utilities such as CopyOnWriteArrayList should not be used.However synchronization can incur a performance overhead. For scenarios where synchronization is not mandatory, then the CopyOnWriteArrayList is a viable, thread-safe alternative to synchronization that leverages multi-core processors and results in higher CPU utilization.
Queue interfaces:
The java.util.Queue interface defines the queue data structure, which stores elements in the order in which they are inserted. New additions go to the end of the line, and elements are removed from the front. It creates a first-in first-out system. This interface is implemented by java.util.LinkedList, java.util.ArrayDeque, and java.util.PriorityQueue.
Queue implementations AbstractQueue class The direct subclasses of AbstractQueue class include ArrayBlockingQueue, ConcurrentLinkedQueue, DelayeQueue, LinkedBlockingDeque, LinkedBlockingQueue.
LinkedTransferQueue and PriorityBlockingQueue.
Note that ArrayDeque and ConcurrentLinkedDeque both extend AbstractCollection but do not extend any other abstact classes such as AbstractQueue.
AbstractQueue is an example of a skeletal implementation.
Queue interfaces:
PriorityQueue class The java.util.PriorityQueue class implements java.util.Queue, but also alters it. PriorityQueue has an additional comparator() method. Instead of elements being ordered in the order in which they are inserted, they are ordered by priority. The method used to determine priority is either the java.lang.Comparable#compareTo(T) method in the elements, or a method given in the constructor. The class creates this by using a heap to keep the items sorted.
Queue interfaces:
ConcurrentLinkedQueue class The java.util.concurrent.ConcurrentLinkedQueue class extends java.util.AbstractQueue. ConcurrentLinkedQueue implements the java.util.Queue interface.The ConcurrentLinkedQueue class is a thread-safe collection, since for any an element placed inside a ConcurrentLinkedQueue, the Java Collection Library guarantees that the element is safely published by allowing any thread to get the element from the collection. An object is said to be safely published if the object's state is made visible to all other thread at the same point in time. Safe publication usually requires synchronization of the publishing and consuming threads.
Queue interfaces:
BlockingQueue interface The java.util.concurrent.BlockingQueue interface extends Queue.The BlockingQueue interface has the following direct sub-interfaces: BlockingDeque and TransferQueue. BlockingQueue works like a regular Queue, but additions to and removals from the BlockingQueue are blocking. If remove(Object o) is called on an empty BlockingQueue, it can be set to wait either a specified time or indefinitely for an item to appear in the BlockingQueue. Similarly, adding an item using the method add(Object o) is subject to an optional capacity restriction on the BlockingQueue, and the method can wait for space to become available in the BlockingQueue before returning. BlockingQueue interface introduces a method take() which removes and gets the head of the BlockingQueue, and waits until the BlockingQueue is no longer empty if required.
Queue interfaces:
Double-ended queue (Deque) interfaces The Deque interface extends the Queue interface. Deque creates a double-ended queue. While a regular Queue only allows insertions at the rear and removals at the front, the Deque allows insertions or removals to take place both at the front and the back. A Deque is like a Queue that can be used forwards or backwards, or both at once. Additionally, both a forwards and a backwards iterator can be generated. The Deque interface is implemented by java.util.ArrayDeque and java.util.LinkedList.
Queue interfaces:
Deque implementations LinkedList class LinkedList, of course, also implements the List interface and can also be used as one. But it also has the Queue methods. LinkedList implements the java.util.Deque interface, giving it more flexibility.
ArrayDeque class ArrayDeque implements the Queue as an array. Similar to LinkedList, ArrayDeque also implements the java.util.Deque interface.
Queue interfaces:
BlockingDeque interface The java.util.concurrent.BlockingDeque interface extends java.util.concurrent.BlockingQueue. BlockingDeque is similar to BlockingQueue. It provides the same methods for insertion and removal with time limits for waiting for the insertion or removal to become possible. However, the interface also provides the flexibility of a Deque. Insertions and removals can take place at both ends. The blocking function is combined with the Deque function.
Set interfaces:
Java's java.util.Setinterface defines the Set. A Set can't have any duplicate elements in it. Additionally, the Set has no set order. As such, elements can't be found by index. Set is implemented by java.util.HashSet, java.util.LinkedHashSet, and java.util.TreeSet.
Set interface implementations There are several implementations of the Set interface, including AbstractSet and its subclasses, and the final static inner class ConcurrentHashMap.KeySetView<K,V> (where K and V are formal type parameters).
AbstractSet AbstractSet is a skeletal implementation for the Set interface.Direct subclasses of AbstractSet include ConcurrentSkipListSet, CopyOnWriteArraySet, EnumSet, HashSet and TreeSet.
Set interfaces:
EnumSet class The EnumSet class extends AbstractSet. The EnumSet class has no public constructors, and only contain static factory methods.EnumSet contains the static factory method EnumSet.of(). This method is an aggregation method. It takes in several parameters, takes into account of the type of the parameters, then returns an instance with the appropriate type. As of 2018, In Java SE8 OpenJDK implementation uses two implementations of EnumSet which are invisible to the client, which are RegularEnumSet and JumboEnumSet. If the RegularEnumSet no longer provided any performance benefits for small enum types, it could be removed from the library without negatively impacting the Java Collection Library.EnumSet is a good replacement for the bit fields, which is a type of set, as described below.Traditionally, whenever developers encountered elements of an enumerated type that needs to be placed in a set, the developer would use the int enum pattern in which every constant is assigned a different power of 2. This bit representation enables the developer to use the bitwise OR operation, so that the constants can be combined into a set, also known as a bit field. This bit field representation enables the developer to make efficient set-based operations and bitwise arithmetic such as intersection and unions.However, there are many problems with bit field representation approach. A bit field is less readable than an int enum constant. Also, if the elements are represented by bit fields, it is impossible to iterate through all of these elements.A recommended alternative approach is to use an EnumSet, where an int enum is used instead of a bit field. This approach uses an EnumSet to represent the set of values that belong to the same Enum type. Since the EnumSet implements the Set interface and no longer requires the use of bit-wise operations, this approach is more type-safe. Furthermore, there are many static factories that allow for object instantiation, such as the method EnumSet.of() method.After the introduction of the EnumSet, the bit field representation approach is considered to be obsolete.
Set interfaces:
HashSet class HashSet uses a hash table. More specifically, it uses a java.util.LinkedHashMap to store the hashes and elements and to prevent duplicates.
LinkedHashSet class The java.util.LinkedHashSet class extends HashSet by creating a doubly linked list that links all of the elements by their insertion order. This ensures that the iteration order over the Set is predictable.
CopyOnWriteArraySet class CopyOnWriteArraySet is a concurrent replacement for a synchronized Set. It provides improved concurrency in many situations by removing the need to perform synchronization or making a copy of the object during iteration, similar to how CopyOnWriteArrayList acts as the concurrent replacement for a synchronized List.
On the other hand, similar to CopyOnWriteArrayList, CopyOnWriteArraySet should not be used when sychronization is mandatory.
Set interfaces:
SortedSet interface The java.util.SortedSet interface extends the java.util.Set interface. Unlike a regular Set, the elements in a SortedSet are sorted, either by the element's compareTo(T o) method, or a method provided to the constructor of the SortedSet. The first and last elements of the SortedSet can be retrieved using the first() and last() methods respectively, and subsets can be created via minimum and maximum values, as well as beginning or ending at the beginning or ending of the SortedSet. The java.util.TreeSet class implements the SortedSet interface.
Set interfaces:
NavigableSet interface The java.util.NavigableSet interface extends the java.util.SortedSet interface and has a few additional methods. The floor(E e), ceiling(E e), lower(E e), and higher(E e) methods find an element in the set that's close to the parameter. Additionally, a descending iterator over the items in the Set is provided. As with SortedSet, java.util.TreeSet implements NavigableSet.
TreeSet class java.util.TreeSet uses a red–black tree implemented by a java.util.TreeMap. The red–black tree ensures that there are no duplicates. Additionally, it allows TreeSet to implement java.util.SortedSet.
ConcurrentSkipListSet class ConcurrentSkipListSet acts as a concurrent replacement for implementations of a synchronized SortedSet. For example it replaces a TreeSet that has been wrapped by the sychronizedMap method.
Map interfaces:
Maps are defined by the java.util.Map interface in Java.
Map interface implementations Maps are data structures that associate a key with an element. This lets the map be very flexible. If the key is the hash code of the element, the Map is essentially a Set. If it's just an increasing number, it becomes a list. Examples of Map implementations include java.util.HashMap, java.util.LinkedHashMap , and java.util.TreeMap.
AbstractMap class AbstractMap is an example of a skeletal implementation.The direct subclasses of AbstractMap class include ConcurrentSkipListMap, EnumMap, HashMap, IdentityHashMap, TreeMap and WeakHashMap.
EnumMap EnumMap extends AbstractMap. EnumMap has comparable speed with an ordinal-indexed array. This is because EnumMap internally uses an array, with implementation details completely hidden from the developer. Hence, the EnumMap gets the type safety of a Map while the performance advantages of an array.
HashMap HashMap uses a hash table. The hashes of the keys are used to find the elements in various buckets. The HashMap is a hash-based collection.
Map interfaces:
LinkedHashMap LinkedHashMap extends HashMap by creating a doubly linked list between the elements, allowing them to be accessed in the order in which they were inserted into the map. LinkedHashMap contains a protected removeEldestEntry method which is called by the put method whenever a new key is added to the Map. The Map removes its eldest entry whenever removeEldestEntry returns true. The removeEldestEntry method can be overridden.
Map interfaces:
TreeMap TreeMap, in contrast to HashMap and LinkedHashMap, uses a red–black tree. The keys are used as the values for the nodes in the tree, and the nodes point to the elements in the Map.
ConcurrentHashMap ConcurrentHashMap is similar to HashMap and is also a hash-based collection. However, there are a number of differences, such as the differences in the locking strategy they use.
The ConcurrentHashMap uses a completely different locking strategy to provide improved scalability and concurrency. ConcurrentHashMap does not synchronize every method using the same lock. Instead, ConcurrentHashMap use a mechanism known as lock striping. This mechanism provides a finer-grained locking mechanism. It also permits a higher degree of shared access.
ConcurrentSkipListMap class ConcurrentSkipListMap acts as a concurrent replacement for implementations of a synchronized SortedMap. ConcurrentSkipListMap is very similar to ConcurrentSkipListSet, since ConcurrentSkipListMap replaces a TreeMap that has been wrapped by the sychronizedMap method.
Map interfaces:
Map subinterfaces SortedMap interface The java.util.SortedMap interface extends the java.util.Map interface. This interface defines a Map that's sorted by the keys provided. Using, once again, the compareTo() method or a method provided in the constructor to the SortedMap, the key-element pairs are sorted by the keys. The first and last keys in the Map can be called by using the firstKey() and lastKey() methods respectively. Additionally, submaps can be created from minimum and maximum keys by using the subMap(K fromKey, K toKey) method. SortedMap is implemented by java.util.TreeMap.
Map interfaces:
NavigableMap interface The java.util.NavigableMap interface extends java.util.SortedMap in various ways. Methods can be called that find the key or map entry that's closest to the given key in either direction. The map can also be reversed, and an iterator in reverse order can be generated from it. It's implemented by java.util.TreeMap.
ConcurrentMap interface The java.util.concurrent.ConcurrentMap interface extends the java.util.Map interface. This interface a thread Safe Map interface, introduced as of Java programming language's Java Collections Framework version 1.5.
Extensions to the Java collections framework:
Java collections framework is extended by the Apache Commons Collections library, which adds collection types such as a bag and bidirectional map, as well as utilities for creating unions and intersections.Google has released its own collections libraries as part of the guava libraries. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**IMG (file format)**
IMG (file format):
IMG, in computing, refers to binary files with the .img filename extension that store raw disk images of floppy disks, hard drives, and optical discs or a bitmap image – .img.
Overview:
The .img filename extension is used by disk image files, which contain raw dumps of a magnetic disk or of an optical disc. Since a raw image consists of a sector-by-sector binary copy of the source medium, the actual format of the file contents will depend on the file system of the disk from which the image was created (such as a version of FAT). Raw disk images of optical media (such as CDs and DVDs) contain a raw image of all the tracks in a disc (which can include audio, data and video tracks). In the case of CD-ROMs and DVDs, these images usually include not only the data from each sector, but the control headers and error correction fields for each sector as well. Since IMG files hold no additional data beyond the disk contents, these files can only be automatically handled by programs that can detect their file systems. For instance, a typical raw disk image of a floppy disk begins with a FAT boot sector, which can be used to identify its file system. Disc images of optical media are usually accompanied by a descriptor file which describes the layout of the disc, and includes information such as track limits which are not stored in the raw image file.
Overview:
Filename extensions and variants The .img file extension was originally used for floppy disk raw disk images only. A similar file extension, .ima, is also used to refer to floppy disk image files by some programs. A variant of IMG, called IMZ, consists of a gzipped version of a raw floppy disk image. These files use the .imz file extension, and are commonly found in compressed images of floppy disks created by WinImage.
Overview:
QEMU uses the .img file extension for raw images of hard drive disks, calling the format simply "raw".
CloneCD stores optical disc images in .img files and generates additional CloneCD Control Files (with .ccd extension) for each image to hold the necessary metadata. The CUE/BIN format stores disc images in .bin files, which are functionally equivalent to .img image files, and uses .cue files as descriptor files.
Overview:
Size The file size of a raw disk image is always a multiple of the sector size. For floppy disks and hard drives this size is typically 512 bytes (but other sizes such as 128 and 1024 exist). More precisely, the file size of a raw disk image of a magnetic disk corresponds to: Cylinders × Heads × (Sectors per track) × (Sector size)E.g. for 80 cylinders (tracks) and 2 heads (sides) with 18 sectors per track: 80 × 2 × 18 × 512 = 1,474,560 bytes or 1440 KBFor optical discs such as CDs and DVDs, the raw sector size is usually 2,352, making the size of a raw disc image a multiple of this value.
Comparison to ISO images:
ISO images are another type of optical disc image files, which commonly use the .iso file extension, but sometimes use the .img file extension as well. They are similar to the raw optical disc images, but contain only one track with computer data obtained from an optical disc. They cannot contain multiple tracks, nor audio or video tracks. They also do not contain the control headers and error correction fields of CD-ROM or DVD sectors that raw disc images usually store. Their internal format follows the structure of an optical disc file system, commonly ISO 9660 (for CDs) or UDF (for DVDs). The CUE/BIN and CCD/IMG formats, which usually contain raw disc images, can also store ISO images instead.
IMG as an image file format:
.img is also a planar bitmap graphics file using simple run-length encoding, originating with Digital Research's GEM. It was commonly used on the Atari ST line of home computers, but also with some GEM-based PC software such as Corel Ventura or Timeworks Publisher.
Other disk image files:
In addition, .img is an Apple Disk Image used by the Mac OS X or macOS operating system.
Garmin .img is a hard-disk image file format which contains a header and many subfiles and used to store the maps for its GPS units.
Tools:
The raw IMG file format is used by several tools: RaWrite and WinImage use the IMG disk image format to read and write floppy disk images.
ImDisk and Virtual Floppy Drive can mount a raw image of a floppy disk to emulate a floppy drive under Microsoft Windows.
Nero Burning ROM supports reading IMG files for creating bootable CDs.
mtools allows manipulation of MS-DOS floppy disk images in Unix systems.
Programs such as dsktrans from the LibDsk suite of command-line tools (available for Linux, MS-DOS, and Microsoft Windows) will convert between different raw disk image formats.
dd can be used in Unix to create raw disk image files of disks.
QEMU uses IMG files as its default format for hard drive disk images.
IrfanView with the plugin "FORMATS" (formats.dll) supports viewing GEM-IMG vector graphics.
Garmin MapSource or GPSMapEdit can be used to read Garmin hard-disk image .img format. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Nvidia NVDEC**
Nvidia NVDEC:
Nvidia NVDEC (formerly known as NVCUVID) is a feature in its graphics cards that performs video decoding, offloading this compute-intensive task from the CPU.It is accompanied by NVENC for video encoding in Nvidia's Video Codec SDK.
Technology:
NVDEC can offload video decoding to full fixed-function decoding hardware (Nvidia PureVideo), or (partially) decode via CUDA software running on the GPU, if fixed-function hardware is not available.Depending on the GPU architecture, the following codecs are supported: MPEG-2 VC-1 H.264 (AVC) H.265 (HEVC) VP8 VP9 AV1
Versions:
NVCUVID was originally distributed as part of the Nvidia CUDA Toolkit. Later, it was renamed to NVDEC and moved to the Nvidia Video Codec SDK.
Operating system support:
NVDEC is available for Windows and Linux operating systems. As NVDEC is a proprietary API (as opposed to the open-source VDPAU API), it is only supported by the proprietary Nvidia driver on Linux.
Application and library support:
Gstreamer has supported NVDEC since 2017.
FFmpeg has supported NVDEC since 2017.
mpv has supported NVDEC since 2017 by the use of FFmpeg.
GPU support:
HW accelerated decode and encode are supported on Nvidia GeForce, Quadro, Tesla, and GRID products with Fermi or newer generation GPUs. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Biodyl**
Biodyl:
Biodyl is a trademark of Merial for a dietary supplement used in animals. It is manufactured in two formulations: a powder for use in an individual animal's drinking water, and an injectable solution. The injectable solution is available by veterinary prescription in some countries and over the counter in others.
Physical and chemical properties:
Biodyl is formulated as a powder to be given in water, and as an injectable solution. The injectable solution is given under the skin, in the muscle, or in a vein, depending on the species of animal. Its intended uses include reducing physiological stress such as due to being transported, and preventing azoturia in performance animals. The manufacturer's own product information describes Biodyl as an "injection solution containing metabolic constituents (adenosine triphosphoric acid or ATP, magnesium and potassium aspartate, sodium selenite and vitamin B 12) for debility, convalescence and myopathies."Composition: Selenium (as sodium selenite) Potassium aspartate hemihydrate Magnesium aspartate tetrahydrate Methyl parahydroxybenzoate Propyl parahydroxybenzoate Excipient
Legal status:
In the United States, Biodyl is not FDA approved, "in that there is not in effect an approval of an application filed with respect to its intended use or uses".
Adverse effects:
The manufacturer however, states that "Biodyl is safe when used as directed. It has been around from the 1950s and adverse reactions have been exceedingly rare over many years of tracking. Less than one animal in over 2 million doses."
Veterinary use:
Implication in polo pony deaths In April 2009, the sudden deaths of 21 polo ponies at Palm Beach International Polo Club in Florida were attributed by a polo team captain to error or tampering in the team's supply of Biodyl. A newspaper in Argentina reported 3 similar deaths of horses at an international competition in Uruguay.In the United States, concerns about a possible manufacturing error or tampering were lost amid a media outcry about the "illegal" use of "illegal" drugs not approved by the FDA, even "banned" by the FDA. In the US, Biodyl is neither an illegal drug nor a banned drug, but it is an unapproved drug. Although Biodyl is a dietary supplement, a type of product that normally is not subject to FDA approval, FDA approval is required to market injectable solutions (except animal vaccines, which are subject to USDA approval). An Associated Press story misreported an October 2008 FDA refusal to permit commercial importation of the solution as a refusal to approve the solution. In fact, Biodyl is not FDA-approved because the manufacturer has never submitted an application for FDA approval. Also, the FDA may permit the importation of unapproved drugs for personal use for pets.However, on April 23 a new concern emerged when a reputable pharmacy in Ocala, Florida disclosed that in compounding a preparation for the polo ponies which may have been intended to substitute for Biodyl, the pharmacy accidentally used an incorrect quantity of one of the ingredients. Compounding of drugs for use in animals is a subject of concern for the FDA. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Putrescine N-hydroxycinnamoyltransferase**
Putrescine N-hydroxycinnamoyltransferase:
In enzymology, a putrescine N-hydroxycinnamoyltransferase (EC 2.3.1.138) is an enzyme that catalyzes the chemical reaction caffeoyl-CoA + putrescine ⇌ CoA + N-caffeoylputrescineThus, the two substrates of this enzyme are caffeoyl-CoA and putrescine, whereas its two products are CoA and N-caffeoylputrescine.
This enzyme belongs to the family of transferases, specifically those acyltransferases transferring groups other than aminoacyl groups. The systematic name of this enzyme class is caffeoyl-CoA:putrescine N-(3,4-dihydroxycinnamoyl)transferase. Other names in common use include caffeoyl-CoA putrescine N-caffeoyl transferase, PHT, putrescine hydroxycinnamoyl transferase, hydroxycinnamoyl-CoA:putrescine hydroxycinnamoyltransferase, and putrescine hydroxycinnamoyltransferase. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Arachnid locomotion**
Arachnid locomotion:
Arachnid locomotion is the various means by which arachnids walk, run, or jump; they make use of more than muscle contraction, employing additional methods like hydraulic compression. Another adaptation seen especially in larger arachnid variants is inclusion of elastic connective tissues.
Hydraulics:
In most arachnids, hydraulic compression acts as the primary means of extension in several of their hinged leg joints, namely the femur–patella joint and tibia–metatarsus joints or second and third leg joints respectively. Instead of blood, hemolymph is used to move nutrients around inside of the arachnid, and has the secondary function of acting as a hydraulic fluid. When compressed by the body of the arachnid, the hemolymph applies compressive force through channels in the limbs that cause them to extend. This motion is then balanced by flexor muscle to retract the leg joints as needed. Due to hydraulics being used for extension, the flexor muscle is able to be significantly larger than would otherwise be possible without impacting size or weight. Measurable core body volume change can occur during periods of higher compression to the legs, as the sinuses of the body contract to achieve pressurization in specific legs. Aside from the normal gait of the arachnid, in some variants, extremely high pressures are used as a means of jumping, propelling rear legs and allowing for much greater and more sudden motion.
Elastics:
In larger variants of arachnids, such as the tarantulas and hairy desert spiders, another mechanism used for locomotion is an elastic sclerite. These sclerites are semi-rigid connectors between leg segments that allow storage and expending of potential energy. This is used as a supplement or in conjunction with the hydraulics normally employed in those joints, allowing for greater weights to be carried, more rapid and sudden movement when combined with the already pronounced flexor muscle acting in those joints, as well as fine motor control with reduced sudden disruption of hemolymph flow. At higher compression of the joint the stiffness of the sclerite has been found to increase significantly, denoting support even outside of normal tension.
Influence on biomimetic design:
Hydraulic locomotion in arachnids has acted as an inspiration for many modern biomimetic concepts in robotics intended for use by or with people, especially in the field of soft robotics. The use of hydraulics in robotic joints is aimed at replacing the more control heavy nature of modern robotics with a more passive system developed in soft actuation. Various forms of actuation and force transmission can be achieved through these inspired designs, including rotation, lifting, and even damping effects. The passive nature of the hydraulic and elastic extensor mechanisms employed have found use in orthotics projects aimed at assisting joints weakened by age or disease.
Fluid secretion:
An additional method used by some arachnids to improve locomotion is to secrete fluids, characterized by a hydrophobic effect, through the pads on the ends of their legs that are in contact with the walking surface. It has been shown that the arachnid is capable of using the adhesive fluid selectively, meaning it can choose to not secrete the fluid in certain circumstances where it would be unwarranted such as in moist conditions. The use of fluids allow the arachnid better traction through improved shear force for both standard locomotion and also sudden movements such as in jumping and leaping.
Challenges in modelling:
Modelling the hydraulic system used by arachnids has been a challenge in the past due to scale and complexity. Simplified models focusing on individual joints and flow channels using modern imaging such as Micro-CT has allowed for mathematical expressions of pressure and flow acting on the joints. Visualizing the flow of hemolymph in small bodies directly has been difficult due to resolution constraints and lack of contrast causing fluid and soft tissue being indistinguishable, but techniques have been employed using a combination of injected microbubbles as tracers in the hemolymph and synchrotron x-ray contrast imaging to track them. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Missile Impact Location System**
Missile Impact Location System:
The Missile Impact Location System or Missile Impact Locating System (MILS) is an ocean acoustic system designed to locate the impact position of test missile nose cones at the ocean's surface and then the position of the cone itself for recovery from the ocean bottom. The systems were installed in the missile test ranges managed by the U.S. Air Force.The systems were first installed in the Eastern Range, at the time the Atlantic Missile Range, and secondly in the Pacific, then known as the Pacific Missile Range. The Atlantic Missile Impact Location System and Pacific Missile Impact Location System were installed from 1958 through 1960. Design and development was by American Telephone and Telegraph Company (AT&T), with its Bell Laboratories research and Western Electric manufacturing elements and was to an extent based on the company's technology and experience developing and deploying the Navy's then classified Sound Surveillance System (SOSUS). Early studies were done at Bell Laboratories' Underwater Systems Development Department examined the problem then the Bell System's other organizations began implementation. The company and Navy assets that had installed the first phase of SOSUS, starting in 1951, were engaged on MILS installation and activation.MILS took several forms and each had a unique configuration based on purpose and local water column and bottom conditions. The target arrays were bottom fixed hydrophones connected by cable to the shore stations. A variant, Sonobuoy MILS (SMILS), was composed of bottom mounted hydrophones augmented by air dropped sonobuoys when in use. The third covered wide ocean areas with fixed hydrophones at distant shore sites was termed broad ocean area (BOA) MILS. All systems exploited the SOFAR channel, also known as the deep sound channel, for long range sound propagation in the ocean.
Target arrays:
The target arrays received the acoustic effect of an object's impact with the ocean surface then by the effect of an explosive charge with location calculated by the difference in arrival times at the hydrophones arranged to form a rough pentagon with a sixth hydrophone at the center. A particular advantage of the pentagon configuration was that a rapid approximate position could be calculated on simple time sequence of the acoustic wave at the hydrophones with detailed analysis producing a more exact location. The effectiveness depended on placement of the hydrophone in the deep sound channel. Since the downrange islands did not offer ocean bottom at that depth in the required configuration a system of suspended hydrophones was used. The difficulty of computing the calibration results for the Atlantic systems led to development of computer programs that became the standard for MILS operational data solutions. The distant placement of the systems revealed the limitations of the existing world geodetic system with various datum systems based on the local geoid, something that would be solved by satellite systems that would develop the means to tie everything together. Target arrays were high accuracy systems usually covering a target area of about 10 nmi (12 mi; 19 km) radius.The Atlantic MILS target arrays were located down range from Cape Canaveral about 700 nmi (810 mi; 1,300 km) at Grand Turk Island, 1,300 nmi (1,500 mi; 2,400 km) at Antigua and 4,400 nmi (5,100 mi; 8,100 km) at Ascension Island.The Pacific Missile Range (PMR), then Navy managed as a complex of ranges, was one of the three national missile ranges. PMR began installation of a Pacific MILS to support Intermediate Range Ballistic Missile (IRBM) tests with impact areas northeast of Hawaii. That system terminated at the Marine Corps Air Station Kaneohe Bay. The IRBM array was operational November 1958. Tests of the Intercontinental Ballistic Missile (ICBM) required MILS monitoring impacts between Midway Island and Wake Island and between Wake Island and Eniwetok. The ICBM range was operational in May 1959 with two target arrays. One was located about 70 nmi (81 mi; 130 km) northeast of Wake and another in the corridor between Wake and Eniwetok. Shore facilities were at Kaneohe and each of the islands.
Broad ocean area (BOA MILS):
This system has less accuracy but extensive coverage area including whole ocean basins. It would cover test vehicles not making the target or other events not directly related to the accuracy tests. Accuracy was improved by pre test calibration by a ship precisely located by a fixed transponder field releasing SOFAR bombs. The BOA hydrophones were located near the deep sound channel axis and were located at Cape Hatteras, Bermuda, Eleuthera (Bahamas), Grand Turk, Puerto Rico, Antigua, Barbados and Ascension. In the Pacific a BOA system was installed to cover the Wake—Eniwetok—Midway impact area.
Broad ocean area (BOA MILS):
Experimental and other uses The BOA MILS sites were involved in events beyond missile testing. Those included both intentional experiments and acoustic incidents in which they were tasked after the fact to examine records. In some experiments MILS was a major participant while in others participation was mainly monitoring and contributing data.
Broad ocean area (BOA MILS):
An example of that monitoring role is the nuclear shot "Sword Fish" in Operation Dominic in which both MILS and SOSUS operated normally simply making recordings and strip charts for a period before the detonation until several hours after. Data has also been provided to support research and support for the International Monitoring System monitoring for nuclear weapons tests. That effort also monitors earthquakes.
Broad ocean area (BOA MILS):
Acoustic propagation research The Kaneohe BOA array, then part of the Pacific Missile Range, was used in the Long Range Acoustic Propagation Project (LRAPP) series of experiments designated Pacific Acoustics Research Kaneohe—Alaska (PARKA). The experiment was required to develop improved models for predicting performance of antisubmarine detection systems and explain the long detection ranges of two to three thousand miles being observed by SOSUS.The Kaneohe shore facility was the operational control center for PARKA I with a hydrophone, bottom sited at 2,070 ft (630.9 m), serving as the secondary receiving site. The main receiving site was the research platform FLIP with hydrophones suspended at 300 ft (91.4 m), 2,500 ft (762.0 m) and 10,800 ft (3,291.8 m). The MILS hydrophones at Midway and the SOSUS array at Point Sur were also used in the experiment.
Broad ocean area (BOA MILS):
Heard Island Feasibility Test The Ascension BOA site had twelve hydrophones in six pairs cabled to the island. All but two pairs were suspended near the deep sound channel. After amplification the signals were fed into a signal processing system.
Broad ocean area (BOA MILS):
Ascension was one of the observing sites for the Heard Island Feasibility Test conducted to observe both the strength and quality of signals traveling at inter-ocean distances and whether those signals were capable of being used in ocean acoustic tomography. A source ship, Cory Chouest, near Heard Island in the Indian Ocean generated signals that were received at Ascension at some 9,200 km (5,700 mi; 5,000 nmi) distance after passing around Africa. Those signals were received as far away as receiving sites and ships on the east and west coasts of North America.
Broad ocean area (BOA MILS):
Vela incident The Ascension array was one of the systems involved in the Vela incident acoustic signal. Three hydrophones correlated acoustic arrivals with the time and estimated location of the double flash detected by the Vela satellite. The detailed study of the Naval Research Laboratory that was based on models from French nuclear testing in the Pacific concluded the acoustic detection was of a near surface nuclear explosion in the vicinity of the Prince Edward Islands.
Sonobuoy MILS (SMILS):
SMILS was exclusively used to support the Navy's fleet ballistic missile programs under the Strategic Systems Project Office with much of the information classified. The range supported the fixed transponder arrays of ten transponders each on a reimbursable basis. The Atlantic range had seven transponder arrays located from 550 nmi (630 mi; 1,020 km) to 4,700 nmi (5,400 mi; 8,700 km) down range.The sonobuoy type impact area used a sonobuoy field, typically four rings 3 nmi (3.5 mi; 5.6 km) apart with outside diameter of 20 nmi (23 mi; 37 km), sowed by aircraft and the reference transponder field for geodetic position. SMILS was not dependent on an island downrange and intended for use in remote ocean areas. The transponders were fixed with the sonobuoy field deployed as needed. The specially equipped aircraft did immediate processing with detailed analysis performed later ashore. A special sonobuoy interrogated the transponder field for position of the sonobuoy pattern to the geodetic referenced transponders and another special sonobuoy established the relative of the sonobuoys within the pattern. Before the sonobuoy deployment a special buoy gathered the data to determine actual sound velocity at various depths at deployment time. Data could be collected by specially modified Navy P-3 aircraft or an Advanced Range Instrumentation Aircraft. The P-3 aircraft, flown from Naval Air Station Patuxent River by Air Test and Evaluation Squadron One, were modified to receive and record more sonobuoys, a special timing system and a monitoring and quick look capability. The sonobuoys were modified standard types, in particular with additional battery life and frequencies. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Kuso Miso Technique**
Kuso Miso Technique:
Kuso Miso Technique (Japanese: くそみそテクニック) is a Japanese gay manga one-shot written and illustrated by Junichi Yamakawa. It was originally published in 1987 in Bara-Komi, a manga supplement of the gay men's magazine Barazoku. The manga depicts a sexual encounter between two men in a public restroom that is complicated by the need of one of the men to relieve himself. Published in Bara-Komi to relative obscurity, Kuso Miso Technique gained notoriety as an Internet meme in the early 2000s after scanned copies of the manga were posted on Japanese imageboards and online forums.
Plot:
Masaki Michishita, a "typical guy" enrolled in preparatory school, is running to a restroom in a public park when he spots a man wearing a jumpsuit sitting on a nearby bench. The man, Takakazu Abe, unzips his jumpsuit and exposes his penis, asking Masaki, "Why don't we do it?" (やらないか, yaranai ka). They proceed to the restroom to have sexual intercourse.
Plot:
Abe performs fellatio on Michishita, upon which Michishita informs him that he needs to urinate. Abe asks Michishita to take the dominant position in anal sex and urinate inside of him, which Michishita obliges. The pair switch positions, and Michishita informs Abe that he now needs to defecate. Abe is annoyed, but nevertheless suggests to a horrified Michishita to defecate while he penetrates him. The ending is left ambiguous, with Michishita commenting through narration that the encounter with Abe "turned out to be a shit show" (くそみそな結果に終わってしまった, kuso miso na kekka ni owatte shimatta).
Characters:
Masaki Michishita (道下 正樹, Michishita Masaki) A prep school student with no prior sexual experience with men, though he often fantasizes about them. He is strongly attracted to Takakazu Abe at first sight, prompting him to think, "Whoa! Hot guy!" (ウホッ! いい男…, Uho! Ii otoko...).Takakazu Abe (阿部 高和, Abe Takakazu) An auto mechanic. In the story, he sits provocatively on a park bench looking for sex. He is well-endowed and an experienced sexual partner.
Publication:
Kuso Miso Technique was originally published in 1987 in Bara-Komi, a manga supplement of the gay men's magazine Barazoku. The series has subsequently been re-published in two anthologies of works by Yamakawa: Uho~tsu!! Ī Otoko-tachi ~ Yamajun Pāfekuto Yamakawa Jun'ichi Komikku (ウホッ！！いい男たち~ヤマジュン・パーフェクト 山川 純一 コミック), published by Fukkan.com in 2003; and Uho~tsu! ! Ī Otoko-tachi 2 Yamajun mi Happyō Sakuhin-shū Yamakawa Jun'ichi (ウホッ!!いい男たち2 ヤマジュン・未発表作品集 山川 純一), published by Booking in 2009.Copyrights for Yamakawa's works were held by Barazoku editor-in-chief Ito Bungaku, who in 2013 transferred the copyright for Kuso Miso Technique to the production company IKD International. In 2018, the copyrights for all of Yamakawa's works, including Kuso Miso Technique, were transferred to the entertainment company Cyzo.
Impact:
Published to relative obscurity in 1987, Kuso Miso Technique gained notoriety as an Internet meme beginning in 2002 after scanned copies of the manga were posted on the imageboards 2channel and Futaba Channel. The memeifcation of the series, particularly the posting of images edited to include Abe and Michishita's faces mid-coitus, trigged the so-called "Yamajun boom" that saw significant public interest in Junichi Yamakawa's manga, with several anthologies of Yamakawa's works selling out the entirety of their print runs. Shift_JIS art of the manga produced on 2channel also contributed to the boom.Lines from the manga such as "uho! Ii otoko", "yaranai ka", and the interjection "uho!" became popular Internet slang on Japanese forums and imageboards, as well as among the Japanese gay community. "Yaranai" was ranked the 16th in 2007's "Net Slang of the Year" in Japan. In 2009, in a survey by Gadget News asking respondents "which manga do you think is the most interesting?", Kuso Miso Technique placed 11th. A variety of official merchandise related to Kuso Miso Technique has been produced, including t-shirts and body pillow covers.A parody song inspired by the series, "Yaranaika", was created by a 2channel user; the song is composed of homoerotic lyrics set to the melody of "Balalaika" by Koharu Kusumi of Morning Musume. The song became popular on video sites such as Nico Nico Douga, where it has been integrated into one of the site's most popular medleys. Yuichiro Nagashima, a kickboxer famous for cosplaying during matches, once entered the ring dressed in Abe's jumpsuit with the parody song as background music.
Adaptations:
Japanese pornography studio Moodyz released a live-action adaptation of Kuso Miso Technique in 2012. Billed as a "pinnacle of futanari anal", the film stars female pornographic actresses Reiko Sawamura and Uta Kohaku as Abe and Michishita, respectively.On April 1, 2023, it was announced that Anime Tokyo will produce a "medium-length" anime film adaptation of Kuso Miso Technique directed by Maki Itō. Itō, a director at Studio Kingyorio, created a live-action fan video adaptation of Kuso Miso Technique while in junior high school, which by 2023 had over 550,000 views on YouTube. The film will reportedly be developed as an all-ages title, and will expand on the story of the original manga and incorporate elements from Yamakawa's other manga. The adaptation will be funded through a crowdfunding campaign slated to commence in Q3 2023. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Separation of concerns**
Separation of concerns:
In computer science, separation of concerns is a design principle for separating a computer program into distinct sections. Each section addresses a separate concern, a set of information that affects the code of a computer program. A concern can be as general as "the details of the hardware for an application", or as specific as "the name of which class to instantiate". A program that embodies SoC well is called a modular program. Modularity, and hence separation of concerns, is achieved by encapsulating information inside a section of code that has a well-defined interface. Encapsulation is a means of information hiding. Layered designs in information systems are another embodiment of separation of concerns (e.g., presentation layer, business logic layer, data access layer, persistence layer).Separation of concerns results in more degrees of freedom for some aspect of the program's design, deployment, or usage. Common among these is increased freedom for simplification and maintenance of code. When concerns are well-separated, there are more opportunities for module upgrade, reuse, and independent development. Hiding the implementation details of modules behind an interface enables improving or modifying a single concern's section of code without having to know the details of other sections and without having to make corresponding changes to those other sections. Modules can also expose different versions of an interface, which increases the freedom to upgrade a complex system in piecemeal fashion without interim loss of functionality.Separation of concerns is a form of abstraction. As with most abstractions, separating concerns means adding additional code interfaces, generally creating more code to be executed. So despite the many benefits of well-separated concerns, there is often an associated execution penalty.
Implementation:
The mechanisms for modular or object-oriented programming that are provided by a programming language are mechanisms that allow developers to provide SoC. For example, object-oriented programming languages such as C#, C++, Delphi, and Java can separate concerns into objects, and architectural design patterns like MVC or MVP can separate presentation and the data-processing (model) from content. Service-oriented design can separate concerns into services. Procedural programming languages such as C and Pascal can separate concerns into procedures or functions. Aspect-oriented programming languages can separate concerns into aspects and objects.
Implementation:
Separation of concerns is an important design principle in many other areas as well, such as urban planning, architecture and information design. The goal is to more effectively understand, design, and manage complex interdependent systems, so that functions can be reused, optimized independently of other functions, and insulated from the potential failure of other functions.
Implementation:
Common examples include separating a space into rooms, so that activity in one room does not affect people in other rooms, and keeping the stove on one circuit and the lights on another, so that overload by the stove does not turn the lights off. The example with rooms shows encapsulation, where information inside one room, such as how messy it is, is not available to the other rooms, except through the interface, which is the door. The example with circuits demonstrates that activity inside one module, which is a circuit with consumers of electricity attached, does not affect activity in a different module, so each module is not concerned with what happens in the other.
Origin:
The term separation of concerns was probably coined by Edsger W. Dijkstra in his 1974 paper "On the role of scientific thought".
Origin:
Let me try to explain to you, what to my taste is characteristic for all intelligent thinking. It is, that one is willing to study in depth an aspect of one's subject matter in isolation for the sake of its own consistency, all the time knowing that one is occupying oneself only with one of the aspects. We know that a program must be correct and we can study it from that viewpoint only; we also know that it should be efficient and we can study its efficiency on another day, so to speak. In another mood we may ask ourselves whether, and if so: why, the program is desirable. But nothing is gained—on the contrary!—by tackling these various aspects simultaneously. It is what I sometimes have called "the separation of concerns", which, even if not perfectly possible, is yet the only available technique for effective ordering of one's thoughts, that I know of. This is what I mean by "focusing one's attention upon some aspect": it does not mean ignoring the other aspects, it is just doing justice to the fact that from this aspect's point of view, the other is irrelevant. It is being one- and multiple-track minded simultaneously.
Origin:
Fifteen years later, it was evident the term separation of concerns was becoming an accepted idea. In 1989, Chris Reade wrote a book titled Elements of Functional Programming that describes separation of concerns: The programmer is having to do several things at the same time, namely, describe what is to be computed; organise the computation sequencing into small steps; organise memory management during the computation.
Origin:
Reade continues to say, Ideally, the programmer should be able to concentrate on the first of the three tasks (describing what is to be computed) without being distracted by the other two, more administrative, tasks. Clearly, administration is important, but by separating it from the main task we are likely to get more reliable results and we can ease the programming problem by automating much of the administration.
Origin:
The separation of concerns has other advantages as well. For example, program proving becomes much more feasible when details of sequencing and memory management are absent from the program. Furthermore, descriptions of what is to be computed should be free of such detailed step-by-step descriptions of how to do it, if they are to be evaluated with different machine architectures. Sequences of small changes to a data object held in a store may be an inappropriate description of how to compute something when a highly parallel machine is being used with thousands of processors distributed throughout the machine and local rather than global storage facilities.
Origin:
Automating the administrative aspects means that the language implementor has to deal with them, but he/she has far more opportunity to make use of very different computation mechanisms with different machine architectures.
Examples:
Internet protocol stack Separation of concerns is crucial to the design of the Internet. In the Internet Protocol Suite, great efforts have been made to separate concerns into well-defined layers. This allows protocol designers to focus on the concerns in one layer, and ignore the other layers. The Application Layer protocol SMTP, for example, is concerned about all the details of conducting an email session over a reliable transport service (usually TCP), but not in the least concerned about how the transport service makes that service reliable. Similarly, TCP is not concerned about the routing of data packets, which is handled at the Internet Layer.
Examples:
HTML, CSS, JavaScript HyperText Markup Language (HTML), Cascading Style Sheets (CSS), and JavaScript (JS) are complementary languages used in the development of web pages and websites. HTML is mainly used for organization of webpage content, CSS is used for definition of content presentation style, and JS defines how the content interacts and behaves with the user. Historically, this was not the case: prior to the introduction of CSS, HTML performed both duties of defining semantics and style.
Examples:
Subject-oriented programming Subject-oriented programming allows separate concerns to be addressed as separate software constructs, each on an equal footing with the others. Each concern provides its own class-structure into which the objects in common are organized, and contributes state and methods to the composite result where they cut across one another. Correspondence rules describe how the classes and methods in the various concerns are related to each other at points where they interact, allowing composite behavior for a method to be derived from several concerns. Multi-dimensional separation of concerns allows the analysis and composition of concerns to be manipulated as a multi-dimensional "matrix" in which each concern provides a dimension in which different points of choice are enumerated, with the cells of the matrix occupied by the appropriate software artifacts.
Examples:
Aspect-oriented programming Aspect-oriented programming allows cross-cutting concerns to be addressed as primary concerns. For example, most programs require some form of security and logging. Security and logging are often secondary concerns, whereas the primary concern is often on accomplishing business goals. However, when designing a program, its security must be built into the design from the beginning instead of being treated as a secondary concern. Applying security afterwards often results in an insufficient security model that leaves too many gaps for future attacks. This may be solved with aspect-oriented programming. For example, an aspect may be written to enforce that calls to a certain API are always logged, or that errors are always logged when an exception is thrown, regardless of whether the program's procedural code handles the exception or propagates it.
Examples:
Levels of analysis in artificial intelligence In cognitive science and artificial intelligence, it is common to refer to David Marr's levels of analysis. At any given time, a researcher may be focusing on (1) what some aspect of intelligence needs to compute, (2) what algorithm it employs, or (3) how that algorithm is implemented in hardware. This separation of concerns is similar to the interface/implementation distinction in software and hardware engineering.
Examples:
Normalized systems In normalized systems separation of concerns is one of the four guiding principles. Adhering to this principle is one of the tools that helps reduce the combinatorial effects that, over time, get introduced in software that is being maintained. In normalized systems separation of concerns is actively supported by the tools.
SoC via partial classes Separation of concerns can be implemented and enforced via partial classes.
SoC via partial classes in Ruby bear_hunting.rb bear_eating.rb bear_hunger.rb | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Upshaw–Schulman syndrome**
Upshaw–Schulman syndrome:
Upshaw–Schulman syndrome (USS) is the recessively inherited form of thrombotic thrombocytopenic purpura (TTP), a rare and complex blood coagulation disease. USS is caused by the absence of the ADAMTS13 protease resulting in the persistence of ultra large von Willebrand factor multimers (ULVWF), causing episodes of acute thrombotic microangiopathy with disseminated multiple small vessel obstructions. These obstructions deprive downstream tissues from blood and oxygen, which can result in tissue damage and death. The presentation of an acute USS episode is variable but usually associated with thrombocytopenia, microangiopathic hemolytic anemia (MAHA) with schistocytes on the peripheral blood smear, fever and signs of ischemic organ damage in the brain, kidney and heart.
Signs and symptoms:
The presentation of TTP is variable. The initial symptoms, which force the patient to medical care, are often the consequence of lower platelet counts like purpura (present in 90% of patients), ecchymosis and hematoma. Patients may also report signs and symptoms as a result of (microangiopathic) hemolytic anemia, such as (dark) beer-brown urine, (mild) jaundice, fatigue and pallor. Cerebral symptoms of various degree are present in many patients, including headache, paresis, speech disorder, visual problems, seizures and disturbance of consciousness up to coma. The symptoms can fluctuate so that they may only be temporarily present but may reappear again later in the TTP episode. Other unspecific symptoms are general malaise, abdominal, joint and muscle pain. Severe manifestations of heart or lung involvements are rare, although affections are not seldom measurable (such as ECG changes).
Cause:
Genetics The ADAMTS13 (a disintegrin and metalloprotease with thrombospondin type 1 motif 13) gene is located on chromosome 9q34 and encoding 29 exons. The ADAMTS13 protease consists of 1427 amino acids and has several protein domains: The signal peptide is thought to have a role in the secretion, folding and stability of the ADAMTS13 protein. It interacts with membrane phospholipids and protein components of the secretory machinery within the cells.
Cause:
The metalloprotease domain contains the active site, which cleaves the unfolded VWF within the A2 domain, between the amino acids Tyrosine1605 and Methionine1606.
Cause:
The disintegrin domain together with the TSP1 repeat, the following cysteine-rich and spacer domains are necessary for substrate recognition, binding and cleavage of the VWF-A2 domain. First, the spacer domain recognizes VWF, increasing the affinity of ADAMTS13 to VWF. Next the disintegrin-like domain engages in a low-affinity binding. Then the metalloprotease domain interacts with the VWF, similar to a three-step molecular zipper.
Cause:
The TSP1 repeats mediate extracellular matrix protein-protein interactions.
The cystein-rich domain is responsible for attachment, for example, to integrins of different cell membranes.
Cause:
The CUB domains take part in protein-protein interactions with VWF domains, which become exposed under shear stress, and are also involved in the binding and cleavage of VWF. Additionally they are involved in the ADAMTS13 protein secretionDisease causing mutations in ADAMTS13, which can be found in all ADAMTS13 protease domains. result predominantly in impaired ADAMTS13 secretion with or without decreased ADAMTS13 protease activity. More than 120 disease causing mutations and numerous single-nucleotide polymorphisms (SNP) are known today. Residual ADAMTS13 activity has been observed with certain mutations and seems to be associated with a later disease-onset. It has been postulated that some SNPs interact with each other and may amplify or reduce overall ADAMTS13 activity.
Cause:
ADAMTS13 function and pathogenesis The ADAMTS protease family contains enzymes that process collagen, cleave inter-cellular matrix, inhibit angiogenesis and blood coagulation. ADAMTS13 belongs to the zinc metalloproteases, and is mainly expressed in liver stellate cells and endothelial cells, but was also found in other cell types, such as platelets, podocytes in the kidney and several brain cells. The only known role of the ADAMTS13 protease is to cleave VWF multimers. The plasma half-life of administered ADAMTS13 in USS patients is around 2–4 days, whereas the protective effects seems to last longer.Usually USS patients have a severely deficient ADAMTS13 activity of <10% of the normal. In this low range there may be residual ADAMTS13 activity, depending on the underlying mutations.In USS severe ADAMTS13 deficiency is often not enough to induce a (first) acute TTP episode. It primarily occurs when an additional (environmental) trigger is present. Recognized triggers are infections (including mild flu-like upper airway infections), pregnancy, heavy alcohol intake or certain drugs. In these situations, VWF is released from its storage organelles, such as Weibel–Palade bodies and granules of platelets. Increased VWF levels in the circulation are leading to a higher demand of ADAMTS13, which is lacking in USS, and can bring forward a TTP episode.
Pathology:
After secretion, ADAMTS13 is either bound to the endothelial surface or free in the blood stream.
Pathology:
The heightened shear stress in small- and microvessels alters the 3D-structure of VWF from the contracted globular form to its linear form. The linear VWF has now its active binding sites exposed, that are important to start blood coagulation. These sites bind platelets and blood vessel lesions by interlinking the stretched VWF with one another – a blood clot is formed.
Pathology:
In its uncut form, (ultra large) VWF's heightened stickiness and interlinking causes spontaneous platelet binding and blood clotting. The linear VWF exposes the A2 domain, so that in the presence of enough ADAMTS13 activity it gets cut to its normal size. VWF in the normal length loses its heightened stickiness and spontaneous crosslinking activity to only form blood clots when needed.
Diagnosis:
A diagnosis of TTP is based on the clinical symptoms with the concomitant presence of thrombocytopenia (platelet count below 100×109/L) and microangiopathic hemolytic anemia with schistocytes on the blood smear, a negative direct antiglobulin test (Coombs test), elevated levels of hemolysis markers (such as total bilirubin, LDH, free hemoglobin, and an unmeasurable haptoglobin), after exclusion of any other apparent cause.USS can present similar to the following diseases, which have to be excluded: fulminant infections, disseminated intravascular coagulation, autoimmune hemolytic anemia, Evans syndrome, the typical and atypical form of hemolytic uremic syndrome, HELLP (hemolysis, elevated liver enzymes, low platelets) syndrome, pre-eclampsia, heparin-induced thrombocytopenia, cancer that is often accompanied with metastasis, kidney injury, antiphospholipid antibody syndrome, and side effects from hematopoietic stem cell transplantation.Of note is that pregnancy associated affections like pre-eclampsia, eclampsia, and HELLP syndrome can overlap in their presentation as pregnancy can trigger TTP episodes.Patients with fulminant infections, disseminated intravascular coagulation, HELLP syndrome, pancreatitis, liver disease, and other active inflammatory conditions may have reduced ADAMTS13 activity, but almost never a relevant severe ADAMTS13 deficiency <10% of the normal.A severe ADAMTS13 deficiency below 5% or <10% of the normal (depending on the definitions) is highly specific for the diagnosis of TTP. ADAMTS13 activity assays are based on the direct or indirect measurement of VWF-cleavage products. Its activity should be measured in blood samples taken before therapy has started, to prevent false high ADAMTS13 activity. If a severe ADAMTS13 deficiency is present an ADAMTS13 inhibitor assay is needed to distinguish between the acquired, autoantibody-mediated and the congenital form of TTP (USS). The presence of antibodies can be tested by ELISA or functional inhibitor assays. The level of ADAMTS13 inhibitor may be fluctuating over the course of disease and depends on free circulatory antibodies, therefore a onetime negative test result does not always exclude the presence of ADAMTS13 inhibitors and thereby an autoimmune origin of TTP. A severe ADAMTS13 deficiency in the absence of an inhibitor, confirmed on a second time point in a healthy episode of a possible USS patient, usually sets the trigger to perform a molecular analysis of the ADAMTS13 gene to confirm a mutation. In unclear cases a plasma infusion trial can be done, showing a USS in the absence of anti-ADAMTS13-antibodies a full recovery of infused plasma-ADAMTS13 activity as well as a plasma half-life of infused ADAMTS13 activity of 2–4 days. A deficiency of ADAMTS13 activity in first-degree relatives is also a very strong indicator for an Upshaw–Schulman syndrome.
Treatment:
The therapy of an acute TTP episode has to be started as early as possible. The standard treatment is the daily replacement of the missing ADAMTS13 protease in form of plasma infusions or in more severe episodes by plasma exchange. In the latter the patients plasma is replaced by donated plasma. The most common sources of ADAMTS13 is platelet-poor fresh frozen plasma (FFP) or solvent-detergent plasma.
Treatment:
The benefit of plasma exchange compared to plasma infusions alone may result from the additional removal of ULVWF. In general both plasma therapies are well tolerated, several mostly minor complications may be observed. The number of infusion/exchange sessions needed to overcome a TTP episode are variable but usually take less than a week in USS. The intensive plasma-therapy is generally stopped when platelet count increases to normal levels and is stable over several days.
Treatment:
Preventive therapy Not all affected patients seem to need a regular preventive plasma infusion therapy, especially as some reach long-term remission without it. Regular plasma infusions are necessary in patients with frequent relapses and in general situations with increased risk to develop an acute episode (as seen above) such as pregnancy. Plasma infusions are given usually every 2–3 weeks to prevent acute episodes of USS, but are often individually adapted.
Outlook:
Several therapy developments for TTP emerged during recent years. Artificially produced ADAMTS13 has been used in mice and testing in humans has been announced. Another drug in development is targeting VWF and its binding sites, thereby reducing VWF-platelet interaction, especially on ULVWF during a TTP episode. Among several (multi-)national data bases a worldwide project has been launched to diagnose USS patients and collect information about them to gain new insights into this rare disease with the goal to optimize patient care.
Epidemiology:
The incidence of acute TTP in adults is around 1.7–4.5 per million and year. These cases are nearly all due to the autoimmune form of TTP, where autoantibodies inhibit ADAMTS13 activity. The prevalence of USS has not yet been determined, but is assumed to constitute less than 5% of all acute TTP cases. The syndrome's inheritance is autosomal recessive, and is more often caused by compound heterozygous than homozygous mutations. The age of onset is variable and can be from neonatal age up to the 5th–6th decade. The risk of relapses differs between affected individuals. Minimization of the burden of disease can be reached by early diagnosis and initiation of prophylaxis if required.
History:
TTP was first recognized as a disease in 1947 and the name was given according to symptoms and underlying pathophysiology which differed from the already known immune thrombocytopenia. Schulman reported a case of TTP in 1960 and Upshaw published a paper in 1978 about relapsing TTP in a patient whom he had followed for 11 years. In his report Upshaw noted the similarities with the reported case by Schulman and hypothesized that the two cases had similar causes – a missing plasma factor.One year later the disease was named Upshaw-Schulmann Syndrome. In 1996 the VWF-cleaving protease was discovered and in the following year found to be the major issue in TTP's pathogenesis. In 2001 the VWF-cleaving protease was identified as ADAMTS13, the gene was mapped to chromosome 9q34, and the first USS-causing mutations were identified. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Six Avoidances (Chinese Music)**
Six Avoidances (Chinese Music):
In Chinese traditional culture, regarding playing Gu-qin, the player should avoid six kinds of occasions: Avoid the time of learning of someone's death; Avoid the time of crying sorrowfully; Avoid the time of being busy with something else; Avoid the time of being angry; Avoid the time of sex; Avoid the time of astonishment. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Tellurite tellurate**
Tellurite tellurate:
A tellurite tellurate is a chemical compound or salt that contains tellurite and tellurate anions (TeO32- and TeO42-). These are mixed anion compounds. Some have third anions.
Naming:
A tellurite tellurate compound may also be called a tellurate tellurite.
Production:
One way to produce a tellurite tellurate compound is by heating oxides together.
Properties:
Monoclinic and orthorhombic dominate crystal structures of the tellurite tellurates. Most compounds are transparent from near ultraviolet to the near infrared. Te-O bonds cause absorption lines in infrared.
Related:
Related to these are the selenate selenites and sulfate sulfites by varying the chalcogen. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**ZK (framework)**
ZK (framework):
ZK is an open-source Ajax Web application framework, written in Java, that enables creation of graphical user interfaces for Web applications with little required programming knowledge.
The core of ZK consists of an Ajax-based event-driven mechanism, over 123 XUL and 83 XHTML-based components, and a mark-up language for designing user interfaces. Programmers design their application pages in feature-rich XUL/XHTML components, and manipulate them upon events triggered by end user's activity. It is similar to the programming model found in desktop GUI-based applications.
ZK (framework):
ZK uses a server-centric approach in which the content synchronization of components and the event pipe-lining between clients and servers are automatically done by the engine, and Ajax plumbing codes are completely transparent to web application developers. Therefore, the end users get the similar engaged interactivity and responsiveness as a desktop application, while programmers' development retains a similar simplicity to that of desktop applications.
ZK (framework):
ZK does not use the standard web request-response mechanism and does not send form fields to the server by making a GET request with query parameters or a POST request. Instead, AJAX requests are sent to the server to update the internal state of each screen widget. At the browser, ZK only downloads a JSON description of the web page and uses a client renderer to turn that into a UI. It's quite efficient and under closer inspection, does not download everything at once. A look at the traffic between client and the server reveals several requests going back and forth between client and browser until the page rendering eventually completes.
ZK (framework):
The optional client-side customization allows the developer to leverage the client-side resources with the so-called server+client fusion, for customization and to reduce the Ajax traffic.
In addition to component-based programming in a manner similar to Swing, ZK supports a mark-up language for rich user interface definition called ZUML.
ZUML is designed for non-programmer developers to design user interfaces intuitively.
ZUML allows developers to meld different markup languages, such as Mozilla XUL language and XHTML, seamlessly into the same page.
ZUML allows developers to embed scripts in pure Java language (interpreted by BeanShell) and use EL expressions to manipulate the components and access data.
Features:
Simply Java. ZK is renowned for its "Ajax without JavaScript" approach, enabling developers to build rich web applications transparently without any knowledge of Ajax and JavaScript.
Responsive design themes along with Bootstrap support HTML 5 and CSS 3 support Over 100 Ajax components offer UI designers a variety of feature rich components to meet the demands of enterprise Ajax applications.
ZUML makes the design of rich user interfaces similar to authoring HTML pages. ZUML is a variant of XUL inheriting all features available to XML, and separates the UI definition from the run-time logic.
Event-driven component-based model (similar to desktop programming models) supporting multiple event driven GUI design patterns.
Support for Model-View-Controller (MVC), Model-View-Presenter and Model-View-ViewModel (MVVM) design patterns Databinding capability via in-ZUML-page annotations that automates CRUD and state synchronization between UI view and the data and POJO.
Application components, such as spreadsheet, pivot table, and calendar.
Embedding scripting support with Java (Beanshell). This is an important benefit as you can use a unified programming language for both user interface and backend programming. Optional support for other serverside Java scripting in other languages such as JavaScript (Rhino), Ruby (JRuby), Python (Jython) and Groovy.
Support for integrating with existing web frameworks via a JSP custom tag library, JSF support, Portlet, and a Servlet Filter. Has support for Spring and Hibernate.
Extensible in that programmers can create new custom UI controls. Groups of controls can be re-used as a macro component.
Extensive charting with ZK Charts
Differences from XUL:
ZK is a server side framework which emits HTML and thus does not depend on client side presence of Gecko making it portable to any browser. ZK takes ZUML (xul and xhtml) serverside pages as input and outputs dhtml for the browser.
Differences from XUL:
ZK processes user interface logic on the server in Java. This increases choice in the scripting engines and application libraries that can be used to create logic. Presentation effects can be delegated to the browser using the Client-side Actions feature to reduce server load for dhtml effects. Running the application logic on the server in a single application container reduces the amount of cross browser and browser version testing when compared to extensively programming the DOM at the browser.
Differences from XUL:
While ZK and XUL have an overlapping subset of components, ZK has its own unique component sets. The component library can be extended with custom components. Macro components can be defined that are composites of components.
ZK provides serverside databinding annotations to declaratively bind UI components to serverside data.
Look and Feel differences.
System requirements:
JRE version 1.5 or later A Web server supporting Servlet 2.3 or later
Prerequisites of programming skills:
Required Basic knowledge of Java or a Java scripting engine language such as Groovy, Rhino (Java JavaScript), JRuby (Java Ruby) or Jython (Java Python) Basic knowledge of HTML and XUL Optional Knowledge of a scripting language to write the glue logic (running on the server): BeanShell (Java interpreter), JavaScript, Groovy, Ruby, Scala and Python Object Oriented Programming Servlet Programming Ajax JavaScript (client-sided) Declarative Databindings JSTL style Expression Language
ZUML:
ZUML (ZK User Interface Markup Language) is a markup language for rich user interfaces definition.
ZUML is designed for non-programmers to design user interfaces efficiently with the ZUML markup ZUML allows developer to meld different markup languages, such as Mozilla XUL language and XHTML, seamlessly into the same page.
ZUML allows developers to embed script in pure Java language (interpreted by BeanShell) and use EL expressions to manipulate the components and access data.
ZUML is supported by ZK.
Official Documentation : ZUML Reference
Client-side technologies:
ZK is a server-centric framework. Technically you don't need to know about the implementation at the client side. It is how ZK Mobile running on Java Mobile VM is done.
Since ZK 5.0, the so-called Server+client Fusion architecture is introduced. Developers are allowed to access the client-side widgets directly if they want to. ZK Client Engine is based on jQuery. Technically you can use jQuery-compliant libraries and widgets.
ZK Add-Ons:
ZK Charts A charting component with APIs for displaying and controlling Charts from server-side ZK Pivottable An Ajax data summarization component ZK Spreadsheet An online Web spreadsheet component. Replaced by Keikai | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Krones**
Krones:
Krones AG is a German packaging and bottling machine manufacturer. It produces lines for filling beverages in plastic and glass bottles or beverage cans. The company manufactures stretch blow-moulding machines for producing polyethylene terephthalate (PET) bottles, plus fillers, labellers, bottle washers, pasteurisers, inspectors, packers and palletisers. This product portfolio is complemented by material flow systems and process technology for producing beverages for breweries, dairies and soft-drink companies.
History:
Krones' corporate evolution is closely connected to the socioeconomic conditions prevalent in Germany following World War II. Hermann Kronseder, the father of the present-day chairman of the supervisory board, used his own designs to manufacture semi-automatic labellers starting with 1951. In the further development as from the 1960s, the firm's range of machinery was extended to include packers and filling systems. In 1980, the company was converted into a stock corporation as Krones AG.
History:
With acquisitions of other companies the present-day range of machines for the beverage industry was reached: 1983: Anton Steinecker Maschinenfabrik (brewhouse manufacture), Freising, Germany 1988: Zierk Maschinenbau GmbH (bottle washers), Flensburg, Germany 1998: Max Kettner GmbH (packaging machines), Rosenheim, Germany 2000: Sander Hansen A/S (pasteurising systems), Brøndby, Denmark 2015: Gernep Group (labellers), Barbing, Germany 2016: System Logistics S.p.A. (60% of shares), Fiorano Modenese, Italy
Corporate data:
The headquarters of the group is situated in Neutraubling near Regensburg, Germany. In Germany in total 10,733 people are employed. Machines and systems are manufactured at the German production facilities (Neutraubling, Nittenau, Flensburg, Freising and Rosenheim). From 2019 on a fabrication site in Debrecen, Hungary, completes production facilities. The internationally focused company achieves more than 90% of its total turnover abroad and is represented worldwide through around 90 subsidiaries and shareholdings. The intralogistics business of Krones is handled by Syskron Holding GmbH since 2014. Subsidiaries KIC Krones GmbH (high-tech adhesives for labels and carton packages, plus processing and operating materials), Neutraubling, Germany KOSME S.R.L. (filling and packaging machines for mid-tier companies), Roverbella, ItalyCover additional market segments in beverage filling.
Corporate data:
In the year 2019 the enterprise held 5,877 patents and utility models.
Corporate structure:
Plastics engineering Stretch blow-moulding machines for the production of PET bottles of up to a volume of 3 liters, with an output of 12,800 to 90,000 bottles/hour. The PET recycling system is based on a PET Flakes cleaning process equipped with progressive temperature controls and decontamination.
Corporate structure:
Filling and packing technology Rinse, filling and capping lines with a rotary concept. With rotary lines, high-speed tasks of up to 72,000 bottles/hour or approximately 120,000 cans/hour are possible, including aseptic filling systems for beverages with a high pH value (> 4,5). For the disinfection of containers and closures, PES or H2O2 is used. Further steps in beverage production as bottle-washing machines, inspection and control systems for bottles and bundles, as well as labelling machines for cold and hot glue or self-adhesive labelling complete the product range. Packaging machines for bundles either one-way or returnable, sorting and grouping stations, as well as palletisation systems supplement the spectrum.
Corporate structure:
Process engineering Breweries can be completely equipped with brewhouses, including fermenting and storage cellar equipment, along with assigned supply installations. For manufacturing plants of non-alcoholic beverages' syrup areas, mixing and carbonising equipment are supplied. Heating systems as UHT- and flash heating systems or pasteurizing systems are available for beverage preservation.
Corporate structure:
IT solutions and material handling systems The control of production process and the integration of production data into an ERP system. Logistics systems provide production and distribution with raw, operating and auxiliary materials, as well as finished products, with either block storage or automatic warehouse systems, including commissioning equipment, forklift guidance and yard management systems for logistics process in beverage and food plants. Activities for logistics have been pooled in System Logistics, a legally independent unit.
Notable milestones since 1997:
1997 Manufacturing of stretch blow-moulding machines for the production of PET bottles 2000 First line for aseptic filling of sensitive soft drinks 2002 First PET recycling line for recovering PET bottles and re-using them as food-grade raw material 2005 Expansion of the firm's aseptic filling technology to include dry sterilisation using H2O2; Krones had already been offering sterilisation with peracetic acid, and is thus the only company to offer both of these technologies.
Notable milestones since 1997:
2009 Manufacturing of water treatment systems 2010 FlexWave heating system for preforms based on microwave technology for energy saving manufacturing of PET bottles.
2011 LitePac: PET bottle formations are provided with strapping tape and a carrying handle only, so that waste from packaging can be reduced by 75% in relation to former shrink film packaging.
2015 DecoType Select: selectively direct printing to grooved and relief structures
Business figures:
Executive board Christoph Klenk, Chairman Uta Anders, CFO Thomas Ricker Markus Tischer Ralf GoldbrunnerChairman of the supervisory board: Volker Kronseder | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Continuous dual Hahn polynomials**
Continuous dual Hahn polynomials:
In mathematics, the continuous dual Hahn polynomials are a family of orthogonal polynomials in the Askey scheme of hypergeometric orthogonal polynomials. They are defined in terms of generalized hypergeometric functions by Sn(x2;a,b,c)=3F2(−n,a+ix,a−ix;a+b,a+c;1).
Roelof Koekoek, Peter A. Lesky, and René F. Swarttouw (2010, 14) give a detailed list of their properties.
Closely related polynomials include the dual Hahn polynomials Rn(x;γ,δ,N), the continuous Hahn polynomials pn(x,a,b, a, b), and the Hahn polynomials. These polynomials all have q-analogs with an extra parameter q, such as the q-Hahn polynomials Qn(x;α,β, N;q), and so on.
Relation to other polynomials:
Wilson polynomials are a generalization of continuous dual Hahn polynomials | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Plastics in the construction industry**
Plastics in the construction industry:
Plastic is the generic name for a family of synthetic materials derived from petrochemicals. It is often product of two or more components. There are many families of plastics and polymers being used in construction industry. Examples of plastics used in building are: Acrylic Composites Expanded Polystyrene Polycarbonate Polyethylene Polypropylene Polyvinyl Chloridein building materials. Some of these properties are:
Merits:
Plastics are strong yet lightweight, and so they are easy to transport & manoeuvre.
They are durable, knock-and scratch resistant with excellent weatherability.
They do not rot or corrode.
Plastics are easy to install; many have a snap-fit kind of jointing procedures.
Plastics offer limitless possibilities in design achieved by extrusion, bending, moulding etc.
They can be given any range of colours by adding pigments.
The plastics are low conductors of heat and thus are used as insulation materials in green building concepts.
The plastics products can achieve tight seals.
Plastic doesn't break easily They can be sawn and nailed employing standard carpentry tools and skills.
They can be easily removed and recycled.
They are poor conductors of electricity.
Disadvantages and limitations:
Plastics may be degraded under the action of direct sunlight which reduces their mechanical strength.
Many plastics are flammable unless treated.
High embodied energy content Low modulus of elasticity: makes them unsuitable for load-bearing applications.
Thermoplastics are subject to creep and soften at moderate temperatures.
Thermal expansion for most plastics is high: adequate thermal movement has to be allowed in detailing.
Many types of plastics are not biodegradable thus cause pollution when they accumulate.
Products:
Some of the examples below are Products of Plastics in the Construction industry: Pipes : Electrical Conduits, Rain Water & Sewage pipes, Plumbing, Gas Distributions.
Cables : PVC Insulation on cables, Insulation Tapes .
Floorings : Flooring tiles & Rolls .
Domes / sky lights : Opaque as well as transparent.
Roofing : Coloured or Double skinned for insulation.
Windows & doors : Extruded sections for Door and windows and panels.
Storage tanks : Storage tanks.
Hardware accessories : Washers, Nut bolts, Sleeves, Anchoring wires.
Temporary structures: Guard cabins, tents Insulation materials: PVC sheets, insulating membranes. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Nasalis muscle**
Nasalis muscle:
The nasalis muscle is a sphincter-like muscle of the nose. It has a transverse part and an alar part. It compresses the nasal cartilages, and can "flare" the nostrils. Some people can use it to close the nostrils to prevent entry of water when underwater. It can be used to test the facial nerve (VII), which supplies it.
Structure:
The nasalis muscle covers the nasal cartilages of the lower surface of the nose. It consists of two parts, transverse and alar: The transverse part (compressor naris muscle) arises from the maxilla, above and lateral to the incisive fossa. Its fibers proceed upward and medially, expanding into a thin aponeurosis which is continuous on the bridge of the nose with that of the muscle of the opposite side, and with the aponeurosis of the procerus muscle. It compresses the nostrils and may completely close them.
Structure:
The alar part (dilator naris muscle) arises from the maxilla over the lateral incisor and inserts into the greater alar cartilage. Its medial fibres tend to blend with the depressor septi nasi muscle, and has been described as part of that muscle.
Nerve supply Like all the other muscles of facial expression, the nasalis muscle is supplied by the facial nerve (VII).
Function:
The nasalis muscle compresses the nasal cartilages. It may also "flare" the nostrils. Some people can use it to close the nostrils to prevent entry of water when underwater.
Clinical significance:
Cleft lip and cleft palate The nasalis muscle is one of the key muscles not formed or inserted correctly with cleft lip and cleft palate deformity. The head of the transverse part needs to be identified during reconstructive surgery so that it can be surgically sutured (connected to) the nasal septum. The origin at the maxilla may also be repositioned for better symmetry.
Clinical significance:
Facial nerve testing Due to it being superficial, the nasalis muscle can be used to test the facial nerve. Specifically, it can be used to test the zygomatic branches. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**HP-GL**
HP-GL:
HP-GL, short for Hewlett-Packard Graphics Language and often written as HPGL, is a printer control language created by Hewlett-Packard (HP). HP-GL was the primary printer control language used by HP plotters. It was introduced with the plotter HP-9872 in 1977 and became a standard for almost all plotters. Hewlett-Packard's printers also usually support HP-GL/2 in addition to PCL.
Design:
The language is formed from a series of two letter codes (mnemonics), followed by optional parameters. For instance an arc can be drawn on a page by sending the string: AA100,100,50; This means Arc Absolute, and the parameters place the center of the arc at absolute coordinates 100,100 on the page, with a starting angle of 50 degrees measured counter-clockwise. A fourth optional parameter (not used here) specifies how far the arc continues, and defaults to 5 degrees.
Design:
When first introduced, HP-GL contained the following commands: Formats: [i]: integer formats between -32767 and 32768. No decimal point.
[d]: decimal format between +/- 127.9999. Optional decimal point.
[c]: ASCII character
Examples:
Typical HP-GL files start with a few setup commands, followed by a long string of graphics commands. The file was in ASCII (text file) format, for instance: The coordinate system was based on the smallest units one of the HP plotters could support, and was set to 25 µm (i.e. 40 units per millimeter, 1016 per inch).
The coordinate space was positive or negative floating point numbers, specifically ±230.
HP-GL/2:
The original HP-GL language did not support definition of line width, as this parameter was determined by the pens loaded into the plotter. With the advent of the first inkjet plotters, line width for the "pens" specified within the HP-GL files had to be set at the printer so it would know what line width to print for each pen, a cumbersome and error-prone process. With Hewlett-Packard Graphics Language/2 aka HP-GL/2, definition of line width was introduced into the language and allowed for elimination of this step. Also, among other improvements a binary file format was defined that allowed for smaller files and shorter file transfer times, and the minimal resolution was reduced.
AGL:
HP-GL is related to AGL (A Graphics Language), an extension of the BASIC programming language. AGL was implemented on Hewlett-Packard minicomputers to simplify controlling a plotter. AGL commands describe the desired graphics plotting function, which the computer relays as several HP-GL instructions to the plotter. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Navier–Stokes equations**
Navier–Stokes equations:
The Navier–Stokes equations ( nav-YAY STOHKS) are partial differential equations which describe the motion of viscous fluid substances, named after French engineer and physicist Claude-Louis Navier and Irish physicist and mathematician George Gabriel Stokes. They were developed over several decades of progressively building the theories, from 1822 (Navier) to 1842-1850 (Stokes).
Navier–Stokes equations:
The Navier–Stokes equations mathematically express momentum balance and conservation of mass for Newtonian fluids. They are sometimes accompanied by an equation of state relating pressure, temperature and density. They arise from applying Isaac Newton's second law to fluid motion, together with the assumption that the stress in the fluid is the sum of a diffusing viscous term (proportional to the gradient of velocity) and a pressure term—hence describing viscous flow. The difference between them and the closely related Euler equations is that Navier–Stokes equations take viscosity into account while the Euler equations model only inviscid flow. As a result, the Navier–Stokes are a parabolic equation and therefore have better analytic properties, at the expense of having less mathematical structure (e.g. they are never completely integrable).
Navier–Stokes equations:
The Navier–Stokes equations are useful because they describe the physics of many phenomena of scientific and engineering interest. They may be used to model the weather, ocean currents, water flow in a pipe and air flow around a wing. The Navier–Stokes equations, in their full and simplified forms, help with the design of aircraft and cars, the study of blood flow, the design of power stations, the analysis of pollution, and many other things. Coupled with Maxwell's equations, they can be used to model and study magnetohydrodynamics.
Navier–Stokes equations:
The Navier–Stokes equations are also of great interest in a purely mathematical sense. Despite their wide range of practical uses, it has not yet been proven whether smooth solutions always exist in three dimensions—i.e., whether they are infinitely differentiable (or even just bounded) at all points in the domain. This is called the Navier–Stokes existence and smoothness problem. The Clay Mathematics Institute has called this one of the seven most important open problems in mathematics and has offered a US$1 million prize for a solution or a counterexample.
Flow velocity:
The solution of the equations is a flow velocity. It is a vector field—to every point in a fluid, at any moment in a time interval, it gives a vector whose direction and magnitude are those of the velocity of the fluid at that point in space and at that moment in time. It is usually studied in three spatial dimensions and one time dimension, although two (spatial) dimensional and steady-state cases are often used as models, and higher-dimensional analogues are studied in both pure and applied mathematics. Once the velocity field is calculated, other quantities of interest such as pressure or temperature may be found using dynamical equations and relations. This is different from what one normally sees in classical mechanics, where solutions are typically trajectories of position of a particle or deflection of a continuum. Studying velocity instead of position makes more sense for a fluid, although for visualization purposes one can compute various trajectories. In particular, the streamlines of a vector field, interpreted as flow velocity, are the paths along which a massless fluid particle would travel. These paths are the integral curves whose derivative at each point is equal to the vector field, and they can represent visually the behavior of the vector field at a point in time.
General continuum equations:
The Navier–Stokes momentum equation can be derived as a particular form of the Cauchy momentum equation, whose general convective form is By setting the Cauchy stress tensor σ {\textstyle {\boldsymbol {\sigma }}} to be the sum of a viscosity term τ {\textstyle {\boldsymbol {\tau }}} (the deviatoric stress) and a pressure term − p I {\textstyle -p\mathbf {I} } (volumetric stress), we arrive at where D D t {\textstyle {\frac {\mathrm {D} }{\mathrm {D} t}}} is the material derivative, defined as ∂ ∂ t + u ⋅ ∇ {\textstyle {\frac {\partial }{\partial t}}+\mathbf {u} \cdot \nabla } , ρ {\textstyle \rho } is the (mass) density, u {\textstyle \mathbf {u} } is the flow velocity, ∇ ⋅ {\textstyle \nabla \cdot \,} is the divergence, p {\textstyle p} is the pressure, t {\textstyle t} is time, τ {\textstyle {\boldsymbol {\tau }}} is the deviatoric stress tensor, which has order 2, g {\textstyle \mathbf {g} } represents body accelerations acting on the continuum, for example gravity, inertial accelerations, electrostatic accelerations, and so on.In this form, it is apparent that in the assumption of an inviscid fluid – no deviatoric stress – Cauchy equations reduce to the Euler equations.
General continuum equations:
Assuming conservation of mass we can use the mass continuity equation (or simply continuity equation), to arrive at the conservation form of the equations of motion. This is often written: where ⊗ {\textstyle \otimes } is the outer product: The left side of the equation describes acceleration, and may be composed of time-dependent and convective components (also the effects of non-inertial coordinates if present). The right side of the equation is in effect a summation of hydrostatic effects, the divergence of deviatoric stress and body forces (such as gravity).
General continuum equations:
All non-relativistic balance equations, such as the Navier–Stokes equations, can be derived by beginning with the Cauchy equations and specifying the stress tensor through a constitutive relation. By expressing the deviatoric (shear) stress tensor in terms of viscosity and the fluid velocity gradient, and assuming constant viscosity, the above Cauchy equations will lead to the Navier–Stokes equations below.
General continuum equations:
Convective acceleration A significant feature of the Cauchy equation and consequently all other continuum equations (including Euler and Navier–Stokes) is the presence of convective acceleration: the effect of acceleration of a flow with respect to space. While individual fluid particles indeed experience time-dependent acceleration, the convective acceleration of the flow field is a spatial effect, one example being fluid speeding up in a nozzle.
Compressible flow:
Remark: here, the deviatoric stress tensor is denoted σ {\textstyle {\boldsymbol {\sigma }}} (instead of τ {\textstyle {\boldsymbol {\tau }}} as it was in the general continuum equations and in the incompressible flow section).
The compressible momentum Navier–Stokes equation results from the following assumptions on the Cauchy stress tensor: the stress is Galilean invariant: it does not depend directly on the flow velocity, but only on spatial derivatives of the flow velocity. So the stress variable is the tensor gradient ∇ u {\textstyle \nabla \mathbf {u} } .
the stress is linear in this variable: σ ( ∇ u ) = C : ( ∇ u ) {\textstyle {\boldsymbol {\sigma }}\left(\nabla \mathbf {u} \right)=\mathbf {C} :\left(\nabla \mathbf {u} \right)} , where C {\textstyle \mathbf {C} } is the fourth-order tensor representing the constant of proportionality, called the viscosity or elasticity tensor, and : is the double-dot product.
Compressible flow:
the fluid is assumed to be isotropic, as with gases and simple liquids, and consequently V {\textstyle \mathbf {V} } is an isotropic tensor; furthermore, since the stress tensor is symmetric, by Helmholtz decomposition it can be expressed in terms of two scalar Lamé parameters, the second viscosity λ {\textstyle \lambda } and the dynamic viscosity μ {\textstyle \mu } , as it is usual in linear elasticity: where I {\textstyle \mathbf {I} } is the identity tensor, ε ( ∇ u ) ≡ 1 2 ∇ u + 1 2 ( ∇ u ) T {\textstyle {\boldsymbol {\varepsilon }}\left(\nabla \mathbf {u} \right)\equiv {\frac {1}{2}}\nabla \mathbf {u} +{\frac {1}{2}}\left(\nabla \mathbf {u} \right)^{T}} is the rate-of-strain tensor and ∇ ⋅ u {\textstyle \nabla \cdot \mathbf {u} } is the divergence (i.e. rate of expansion) of the flow. So this decomposition can be explicitly defined as: Since the trace of the rate-of-strain tensor in three dimensions is: The trace of the stress tensor in three dimensions becomes: So by alternatively decomposing the stress tensor into isotropic and deviatoric parts, as usual in fluid dynamics: Introducing the bulk viscosity ζ {\textstyle \zeta } , we arrive to the linear constitutive equation in the form usually employed in thermal hydraulics: Both second viscosity ζ {\textstyle \zeta } and dynamic viscosity μ {\textstyle \mu } need not be constant – in general, they depend on two thermodynamics variables if the fluid contains a single chemical species, say for example, pressure and temperature. Any equation that makes explicit one of these transport coefficient in the conservation variables is called an equation of state.The most general of the Navier–Stokes equations become Apart from its dependence of pressure and temperature, the second viscosity coefficient also depends on the process, that is to say, the second viscosity coefficient is not just a material property. For instance, in the case of a sound wave with a definitive frequency that alternatively compresses and expands a fluid element, the second viscosity coefficient depends on the frequency of the wave. This dependence is called the dispersion. In some cases, the second viscosity ζ {\textstyle \zeta } can be assumed to be constant in which case, the effect of the volume viscosity ζ {\textstyle \zeta } is that the mechanical pressure is not equivalent to the thermodynamic pressure: as demonstrated below.
Compressible flow:
However, this difference is usually neglected most of the time (that is whenever we are not dealing with processes such as sound absorption and attenuation of shock waves, where second viscosity coefficient becomes important) by explicitly assuming ζ = 0 {\textstyle \zeta =0} . The assumption of setting ζ = 0 {\textstyle \zeta =0} is called as the Stokes hypothesis. The validity of Stokes hypothesis can be demonstrated for monoatomic gas both experimentally and from the kinetic theory,; for other gases and liquids, Stokes hypothesis is generally incorrect. With the Stokes hypothesis, the Navier–Stokes equations become If the dynamic viscosity μ is also assumed to be constant, the equations can be simplified further. By computing the divergence of the stress tensor, since the divergence of tensor ∇ u {\textstyle \nabla \mathbf {u} } is ∇ 2 u {\textstyle \nabla ^{2}\mathbf {u} } and the divergence of tensor ( ∇ u ) T {\textstyle \left(\nabla \mathbf {u} \right)^{\mathrm {T} }} is ∇ ( ∇ ⋅ u ) {\textstyle \nabla \left(\nabla \cdot \mathbf {u} \right)} , one finally arrives to the compressible (most general) Navier–Stokes momentum equation: where D D t {\textstyle {\frac {\mathrm {D} }{\mathrm {D} t}}} is the material derivative. The left-hand side changes in the conservation form of the Navier–Stokes momentum equation: Bulk viscosity is assumed to be constant, otherwise it should not be taken out of the last derivative. The convective acceleration term can also be written as where the vector ( ∇ × u ) × u {\textstyle (\nabla \times \mathbf {u} )\times \mathbf {u} } is known as the Lamb vector.
Compressible flow:
For the special case of an incompressible flow, the pressure constrains the flow so that the volume of fluid elements is constant: isochoric flow resulting in a solenoidal velocity field with ∇ ⋅ u = 0 {\textstyle \nabla \cdot \mathbf {u} =0} .
Incompressible flow:
The incompressible momentum Navier–Stokes equation results from the following assumptions on the Cauchy stress tensor: the stress is Galilean invariant: it does not depend directly on the flow velocity, but only on spatial derivatives of the flow velocity. So the stress variable is the tensor gradient {\textstyle \nabla \mathbf {u} } the fluid is assumed to be isotropic, as with gases and simple liquids, and consequently {\textstyle {\boldsymbol {\tau }}} is an isotropic tensor; furthermore, since the deviatoric stress tensor can be expressed in terms of the dynamic viscosity {\textstyle \mu } where is the rate-of-strain tensor. So this decomposition can be made explicit as: Dynamic viscosity μ need not be constant – in incompressible flows it can depend on density and on pressure. Any equation that makes explicit one of these transport coefficient in the conservative variables is called an equation of state.The divergence of the deviatoric stress is given by: because {\textstyle \nabla \cdot \mathbf {u} =0} for an incompressible fluid.
Incompressible flow:
Incompressibility rules out density and pressure waves like sound or shock waves, so this simplification is not useful if these phenomena are of interest. The incompressible flow assumption typically holds well with all fluids at low Mach numbers (say up to about Mach 0.3), such as for modelling air winds at normal temperatures. the incompressible Navier–Stokes equations are best visualized by dividing for the density: If the density is constant throughout the fluid domain, or, in other words, if all fluid elements have the same density, {\textstyle \rho =\rho _{0}} , then we have where {\textstyle \nu ={\frac {\mu }{\rho _{0}}}} is called the kinematic viscosity.
Incompressible flow:
It is well worth observing the meaning of each term (compare to the Cauchy momentum equation): The higher-order term, namely the shear stress divergence {\textstyle \nabla \cdot {\boldsymbol {\tau }}} , has simply reduced to the vector Laplacian term {\textstyle \mu \nabla ^{2}\mathbf {u} } . This Laplacian term can be interpreted as the difference between the velocity at a point and the mean velocity in a small surrounding volume. This implies that – for a Newtonian fluid – viscosity operates as a diffusion of momentum, in much the same way as the heat conduction. In fact neglecting the convection term, incompressible Navier–Stokes equations lead to a vector diffusion equation (namely Stokes equations), but in general the convection term is present, so incompressible Navier–Stokes equations belong to the class of convection–diffusion equations.
Incompressible flow:
In the usual case of an external field being a conservative field: by defining the hydraulic head: one can finally condense the whole source in one term, arriving to the incompressible Navier–Stokes equation with conservative external field: The incompressible Navier–Stokes equations with conservative external field is the fundamental equation of hydraulics. The domain for these equations is commonly a 3 or less dimensional Euclidean space, for which an orthogonal coordinate reference frame is usually set to explicit the system of scalar partial differential equations to be solved. In 3-dimensional orthogonal coordinate systems are 3: Cartesian, cylindrical, and spherical. Expressing the Navier–Stokes vector equation in Cartesian coordinates is quite straightforward and not much influenced by the number of dimensions of the euclidean space employed, and this is the case also for the first-order terms (like the variation and convection ones) also in non-cartesian orthogonal coordinate systems. But for the higher order terms (the two coming from the divergence of the deviatoric stress that distinguish Navier–Stokes equations from Euler equations) some tensor calculus is required for deducing an expression in non-cartesian orthogonal coordinate systems.
Incompressible flow:
The incompressible Navier–Stokes equation is composite, the sum of two orthogonal equations, where {\textstyle \Pi ^{S}} and {\textstyle \Pi ^{I}} are solenoidal and irrotational projection operators satisfying {\textstyle \Pi ^{S}+\Pi ^{I}-1} and {\textstyle \mathbf {f} ^{S}} and {\textstyle \mathbf {f} ^{I}} are the non-conservative and conservative parts of the body force. This result follows from the Helmholtz theorem (also known as the fundamental theorem of vector calculus). The first equation is a pressureless governing equation for the velocity, while the second equation for the pressure is a functional of the velocity and is related to the pressure Poisson equation.
Incompressible flow:
The explicit functional form of the projection operator in 3D is found from the Helmholtz Theorem: with a similar structure in 2D. Thus the governing equation is an integro-differential equation similar to Coulomb and Biot–Savart law, not convenient for numerical computation.
Incompressible flow:
An equivalent weak or variational form of the equation, proved to produce the same velocity solution as the Navier–Stokes equation, is given by, for divergence-free test functions {\textstyle \mathbf {w} } satisfying appropriate boundary conditions. Here, the projections are accomplished by the orthogonality of the solenoidal and irrotational function spaces. The discrete form of this is eminently suited to finite element computation of divergence-free flow, as we shall see in the next section. There one will be able to address the question "How does one specify pressure-driven (Poiseuille) problems with a pressureless governing equation?".
Incompressible flow:
The absence of pressure forces from the governing velocity equation demonstrates that the equation is not a dynamic one, but rather a kinematic equation where the divergence-free condition serves the role of a conservation equation. This all would seem to refute the frequent statements that the incompressible pressure enforces the divergence-free condition.
Incompressible flow:
Weak form of the incompressible Navier–Stokes equations Strong form Consider the incompressible Navier–Stokes equations for a Newtonian fluid of constant density {\textstyle \rho } in a domain with boundary being {\textstyle \Gamma _{D}} and {\textstyle \Gamma _{N}} portions of the boundary where respectively a Dirichlet and a Neumann boundary condition is applied ( {\textstyle \Gamma _{D}\cap \Gamma _{N}=\emptyset } ): {\textstyle \mathbf {u} } is the fluid velocity, {\textstyle p} the fluid pressure, {\textstyle \mathbf {f} } a given forcing term, n^ the outward directed unit normal vector to {\textstyle \Gamma _{N}} , and {\textstyle {\boldsymbol {\sigma }}(\mathbf {u} ,p)} the viscous stress tensor defined as: Let {\textstyle \mu } be the dynamic viscosity of the fluid, {\textstyle \mathbf {I} } the second-order identity tensor and {\textstyle {\boldsymbol {\varepsilon }}(\mathbf {u} )} the strain-rate tensor defined as: The functions {\textstyle \mathbf {g} } and {\textstyle \mathbf {h} } are given Dirichlet and Neumann boundary data, while {\textstyle \mathbf {u} _{0}} is the initial condition. The first equation is the momentum balance equation, while the second represents the mass conservation, namely the continuity equation. Assuming constant dynamic viscosity, using the vectorial identity and exploiting mass conservation, the divergence of the total stress tensor in the momentum equation can also be expressed as: Moreover, note that the Neumann boundary conditions can be rearranged as: Weak form In order to find the weak form of the Navier–Stokes equations, firstly, consider the momentum equation multiply it for a test function {\textstyle \mathbf {v} } , defined in a suitable space {\textstyle V} , and integrate both members with respect to the domain {\textstyle \Omega } Counter-integrating by parts the diffusive and the pressure terms and by using the Gauss' theorem: Using these relations, one gets: In the same fashion, the continuity equation is multiplied for a test function q belonging to a space {\textstyle Q} and integrated in the domain {\textstyle \Omega } The space functions are chosen as follows: Considering that the test function v vanishes on the Dirichlet boundary and considering the Neumann condition, the integral on the boundary can be rearranged as: Having this in mind, the weak formulation of the Navier–Stokes equations is expressed as: Discrete velocity With partitioning of the problem domain and defining basis functions on the partitioned domain, the discrete form of the governing equation is It is desirable to choose basis functions that reflect the essential feature of incompressible flow – the elements must be divergence-free. While the velocity is the variable of interest, the existence of the stream function or vector potential is necessary by the Helmholtz theorem. Further, to determine fluid flow in the absence of a pressure gradient, one can specify the difference of stream function values across a 2D channel, or the line integral of the tangential component of the vector potential around the channel in 3D, the flow being given by Stokes' theorem. Discussion will be restricted to 2D in the following.
Incompressible flow:
We further restrict discussion to continuous Hermite finite elements which have at least first-derivative degrees-of-freedom. With this, one can draw a large number of candidate triangular and rectangular elements from the plate-bending literature. These elements have derivatives as components of the gradient. In 2D, the gradient and curl of a scalar are clearly orthogonal, given by the expressions, Adopting continuous plate-bending elements, interchanging the derivative degrees-of-freedom and changing the sign of the appropriate one gives many families of stream function elements.
Incompressible flow:
Taking the curl of the scalar stream function elements gives divergence-free velocity elements. The requirement that the stream function elements be continuous assures that the normal component of the velocity is continuous across element interfaces, all that is necessary for vanishing divergence on these interfaces.
Boundary conditions are simple to apply. The stream function is constant on no-flow surfaces, with no-slip velocity conditions on surfaces.
Stream function differences across open channels determine the flow. No boundary conditions are necessary on open boundaries, though consistent values may be used with some problems. These are all Dirichlet conditions.
The algebraic equations to be solved are simple to set up, but of course are non-linear, requiring iteration of the linearized equations.
Similar considerations apply to three-dimensions, but extension from 2D is not immediate because of the vector nature of the potential, and there exists no simple relation between the gradient and the curl as was the case in 2D.
Incompressible flow:
Pressure recovery Recovering pressure from the velocity field is easy. The discrete weak equation for the pressure gradient is, where the test/weight functions are irrotational. Any conforming scalar finite element may be used. However, the pressure gradient field may also be of interest. In this case, one can use scalar Hermite elements for the pressure. For the test/weight functions {\textstyle \mathbf {g} _{i}} one would choose the irrotational vector elements obtained from the gradient of the pressure element.
Non-inertial frame of reference:
The rotating frame of reference introduces some interesting pseudo-forces into the equations through the material derivative term. Consider a stationary inertial frame of reference K {\textstyle K} , and a non-inertial frame of reference K ′ {\textstyle K'} , which is translating with velocity U ( t ) {\textstyle \mathbf {U} (t)} and rotating with angular velocity Ω ( t ) {\textstyle \Omega (t)} with respect to the stationary frame. The Navier–Stokes equation observed from the non-inertial frame then becomes Here x {\textstyle \mathbf {x} } and u {\textstyle \mathbf {u} } are measured in the non-inertial frame. The first term in the parenthesis represents Coriolis acceleration, the second term is due to centrifugal acceleration, the third is due to the linear acceleration of K ′ {\textstyle K'} with respect to K {\textstyle K} and the fourth term is due to the angular acceleration of K ′ {\textstyle K'} with respect to K {\textstyle K} .
Other equations:
The Navier–Stokes equations are strictly a statement of the balance of momentum. To fully describe fluid flow, more information is needed, how much depending on the assumptions made. This additional information may include boundary data (no-slip, capillary surface, etc.), conservation of mass, balance of energy, and/or an equation of state.
Other equations:
Continuity equation for incompressible fluid Regardless of the flow assumptions, a statement of the conservation of mass is generally necessary. This is achieved through the mass continuity equation, given in its most general form as: or, using the substantive derivative: For incompressible fluid, density along the line of flow remains constant over time, Therefore divergence of velocity is always zero:
Stream function for incompressible 2D fluid:
Taking the curl of the incompressible Navier–Stokes equation results in the elimination of pressure. This is especially easy to see if 2D Cartesian flow is assumed (like in the degenerate 3D case with u z = 0 {\textstyle u_{z}=0} and no dependence of anything on z {\textstyle z} ), where the equations reduce to: Differentiating the first with respect to y {\textstyle y} , the second with respect to x {\textstyle x} and subtracting the resulting equations will eliminate pressure and any conservative force. For incompressible flow, defining the stream function ψ {\textstyle \psi } through results in mass continuity being unconditionally satisfied (given the stream function is continuous), and then incompressible Newtonian 2D momentum and mass conservation condense into one equation: where ∇ 4 {\textstyle \nabla ^{4}} is the 2D biharmonic operator and ν {\textstyle \nu } is the kinematic viscosity, ν = μ p {\textstyle \nu ={\frac {\mu }{p}}} . We can also express this compactly using the Jacobian determinant: This single equation together with appropriate boundary conditions describes 2D fluid flow, taking only kinematic viscosity as a parameter. Note that the equation for creeping flow results when the left side is assumed zero.
Stream function for incompressible 2D fluid:
In axisymmetric flow another stream function formulation, called the Stokes stream function, can be used to describe the velocity components of an incompressible flow with one scalar function.
Stream function for incompressible 2D fluid:
The incompressible Navier–Stokes equation is a differential algebraic equation, having the inconvenient feature that there is no explicit mechanism for advancing the pressure in time. Consequently, much effort has been expended to eliminate the pressure from all or part of the computational process. The stream function formulation eliminates the pressure but only in two dimensions and at the expense of introducing higher derivatives and elimination of the velocity, which is the primary variable of interest.
Properties:
Nonlinearity The Navier–Stokes equations are nonlinear partial differential equations in the general case and so remain in almost every real situation. In some cases, such as one-dimensional flow and Stokes flow (or creeping flow), the equations can be simplified to linear equations. The nonlinearity makes most problems difficult or impossible to solve and is the main contributor to the turbulence that the equations model.
Properties:
The nonlinearity is due to convective acceleration, which is an acceleration associated with the change in velocity over position. Hence, any convective flow, whether turbulent or not, will involve nonlinearity. An example of convective but laminar (nonturbulent) flow would be the passage of a viscous fluid (for example, oil) through a small converging nozzle. Such flows, whether exactly solvable or not, can often be thoroughly studied and understood.
Properties:
Turbulence Turbulence is the time-dependent chaotic behaviour seen in many fluid flows. It is generally believed that it is due to the inertia of the fluid as a whole: the culmination of time-dependent and convective acceleration; hence flows where inertial effects are small tend to be laminar (the Reynolds number quantifies how much the flow is affected by inertia). It is believed, though not known with certainty, that the Navier–Stokes equations describe turbulence properly.The numerical solution of the Navier–Stokes equations for turbulent flow is extremely difficult, and due to the significantly different mixing-length scales that are involved in turbulent flow, the stable solution of this requires such a fine mesh resolution that the computational time becomes significantly infeasible for calculation or direct numerical simulation. Attempts to solve turbulent flow using a laminar solver typically result in a time-unsteady solution, which fails to converge appropriately. To counter this, time-averaged equations such as the Reynolds-averaged Navier–Stokes equations (RANS), supplemented with turbulence models, are used in practical computational fluid dynamics (CFD) applications when modeling turbulent flows. Some models include the Spalart–Allmaras, k–ω, k–ε, and SST models, which add a variety of additional equations to bring closure to the RANS equations. Large eddy simulation (LES) can also be used to solve these equations numerically. This approach is computationally more expensive—in time and in computer memory—than RANS, but produces better results because it explicitly resolves the larger turbulent scales.
Properties:
Applicability Together with supplemental equations (for example, conservation of mass) and well-formulated boundary conditions, the Navier–Stokes equations seem to model fluid motion accurately; even turbulent flows seem (on average) to agree with real world observations.
Properties:
The Navier–Stokes equations assume that the fluid being studied is a continuum (it is infinitely divisible and not composed of particles such as atoms or molecules), and is not moving at relativistic velocities. At very small scales or under extreme conditions, real fluids made out of discrete molecules will produce results different from the continuous fluids modeled by the Navier–Stokes equations. For example, capillarity of internal layers in fluids appears for flow with high gradients. For large Knudsen number of the problem, the Boltzmann equation may be a suitable replacement. Failing that, one may have to resort to molecular dynamics or various hybrid methods.Another limitation is simply the complicated nature of the equations. Time-tested formulations exist for common fluid families, but the application of the Navier–Stokes equations to less common families tends to result in very complicated formulations and often to open research problems. For this reason, these equations are usually written for Newtonian fluids where the viscosity model is linear; truly general models for the flow of other kinds of fluids (such as blood) do not exist.
Application to specific problems:
The Navier–Stokes equations, even when written explicitly for specific fluids, are rather generic in nature and their proper application to specific problems can be very diverse. This is partly because there is an enormous variety of problems that may be modeled, ranging from as simple as the distribution of static pressure to as complicated as multiphase flow driven by surface tension.
Application to specific problems:
Generally, application to specific problems begins with some flow assumptions and initial/boundary condition formulation, this may be followed by scale analysis to further simplify the problem.
Parallel flow Assume steady, parallel, one-dimensional, non-convective pressure-driven flow between parallel plates, the resulting scaled (dimensionless) boundary value problem is: The boundary condition is the no slip condition. This problem is easily solved for the flow field: From this point onward, more quantities of interest can be easily obtained, such as viscous drag force or net flow rate.
Application to specific problems:
Radial flow Difficulties may arise when the problem becomes slightly more complicated. A seemingly modest twist on the parallel flow above would be the radial flow between parallel plates; this involves convection and thus non-linearity. The velocity field may be represented by a function f(z) that must satisfy: This ordinary differential equation is what is obtained when the Navier–Stokes equations are written and the flow assumptions applied (additionally, the pressure gradient is solved for). The nonlinear term makes this a very difficult problem to solve analytically (a lengthy implicit solution may be found which involves elliptic integrals and roots of cubic polynomials). Issues with the actual existence of solutions arise for R > 1.41 {\textstyle R>1.41} (approximately; this is not √2), the parameter R {\textstyle R} being the Reynolds number with appropriately chosen scales. This is an example of flow assumptions losing their applicability, and an example of the difficulty in "high" Reynolds number flows.
Application to specific problems:
Convection A type of natural convection that can be described by the Navier–Stokes equation is the Rayleigh–Bénard convection. It is one of the most commonly studied convection phenomena because of its analytical and experimental accessibility.
Exact solutions of the Navier–Stokes equations:
Some exact solutions to the Navier–Stokes equations exist. Examples of degenerate cases—with the non-linear terms in the Navier–Stokes equations equal to zero—are Poiseuille flow, Couette flow and the oscillatory Stokes boundary layer. But also, more interesting examples, solutions to the full non-linear equations, exist, such as Jeffery–Hamel flow, Von Kármán swirling flow, stagnation point flow, Landau–Squire jet, and Taylor–Green vortex.
Exact solutions of the Navier–Stokes equations:
Note that the existence of these exact solutions does not imply they are stable: turbulence may develop at higher Reynolds numbers.
Under additional assumptions, the component parts can be separated.
Exact solutions of the Navier–Stokes equations:
A three-dimensional steady-state vortex solution A steady-state example with no singularities comes from considering the flow along the lines of a Hopf fibration. Let {\textstyle r} be a constant radius of the inner coil. One set of solutions is given by: for arbitrary constants {\textstyle A} and {\textstyle B} . This is a solution in a non-viscous gas (compressible fluid) whose density, velocities and pressure goes to zero far from the origin. (Note this is not a solution to the Clay Millennium problem because that refers to incompressible fluids where {\textstyle \rho } is a constant, and neither does it deal with the uniqueness of the Navier–Stokes equations with respect to any turbulence properties.) It is also worth pointing out that the components of the velocity vector are exactly those from the Pythagorean quadruple parametrization. Other choices of density and pressure are possible with the same velocity field: Viscous three-dimensional periodic solutions Two examples of periodic fully-three-dimensional viscous solutions are described in.
Exact solutions of the Navier–Stokes equations:
These solutions are defined on a three-dimensional torus T3=[0,L]3 and are characterized by positive and negative helicity respectively.
The solution with positive helicity is given by: where k=2π/L is the wave number and the velocity components are normalized so that the average kinetic energy per unit of mass is U02/2 at t=0 The pressure field is obtained from the velocity field as p=p0−ρ0‖u‖2/2 (where p0 and ρ0 are reference values for the pressure and density fields respectively).
Since both the solutions belong to the class of Beltrami flow, the vorticity field is parallel to the velocity and, for the case with positive helicity, is given by ω=3ku . These solutions can be regarded as a generalization in three dimensions of the classic two-dimensional Taylor-Green Taylor–Green vortex.
Wyld diagrams:
Wyld diagrams are bookkeeping graphs that correspond to the Navier–Stokes equations via a perturbation expansion of the fundamental continuum mechanics. Similar to the Feynman diagrams in quantum field theory, these diagrams are an extension of Keldysh's technique for nonequilibrium processes in fluid dynamics. In other words, these diagrams assign graphs to the (often) turbulent phenomena in turbulent fluids by allowing correlated and interacting fluid particles to obey stochastic processes associated to pseudo-random functions in probability distributions.
Representations in 3D:
Note that the formulas in this section make use of the single-line notation for partial derivatives, where, e.g. ∂ x u {\textstyle \partial _{x}u} means the partial derivative of u {\textstyle u} with respect to x {\textstyle x} , and ∂ y 2 f θ {\textstyle \partial _{y}^{2}f_{\theta }} means the second-order partial derivative of f θ {\textstyle f_{\theta }} with respect to y {\textstyle y} .
Representations in 3D:
A 2022 paper provides a less costly, dynamical and recurrent solution of the Navier-Stokes equation for 3D turbulent fluid flows. On suitably short time scales, the dynamics of turbulence is deterministic.
Representations in 3D:
Cartesian coordinates From the general form of the Navier–Stokes, with the velocity vector expanded as u = ( u x , u y , u z ) {\textstyle \mathbf {u} =(u_{x},u_{y},u_{z})} , sometimes respectively named u {\textstyle u} , v {\textstyle v} , w {\textstyle w} , we may write the vector equation explicitly, Note that gravity has been accounted for as a body force, and the values of g x {\textstyle g_{x}} , g y {\textstyle g_{y}} , g z {\textstyle g_{z}} will depend on the orientation of gravity with respect to the chosen set of coordinates.
Representations in 3D:
The continuity equation reads: When the flow is incompressible, ρ {\textstyle \rho } does not change for any fluid particle, and its material derivative vanishes: D ρ D t = 0 {\textstyle {\frac {\mathrm {D} \rho }{\mathrm {D} t}}=0} . The continuity equation is reduced to: Thus, for the incompressible version of the Navier–Stokes equation the second part of the viscous terms fall away (see Incompressible flow).
Representations in 3D:
This system of four equations comprises the most commonly used and studied form. Though comparatively more compact than other representations, this is still a nonlinear system of partial differential equations for which solutions are difficult to obtain.
Representations in 3D:
Cylindrical coordinates A change of variables on the Cartesian equations will yield the following momentum equations for r {\textstyle r} , ϕ {\textstyle \phi } , and z {\textstyle z} The gravity components will generally not be constants, however for most applications either the coordinates are chosen so that the gravity components are constant or else it is assumed that gravity is counteracted by a pressure field (for example, flow in horizontal pipe is treated normally without gravity and without a vertical pressure gradient). The continuity equation is: This cylindrical representation of the incompressible Navier–Stokes equations is the second most commonly seen (the first being Cartesian above). Cylindrical coordinates are chosen to take advantage of symmetry, so that a velocity component can disappear. A very common case is axisymmetric flow with the assumption of no tangential velocity ( u ϕ = 0 {\textstyle u_{\phi }=0} ), and the remaining quantities are independent of ϕ {\textstyle \phi } : Spherical coordinates |In spherical coordinates, the r {\textstyle r} , ϕ {\textstyle \phi } , and θ {\textstyle \theta } momentum equations are (note the convention used: θ {\textstyle \theta } is polar angle, or colatitude, 0 ≤ θ ≤ π {\textstyle 0\leq \theta \leq \pi } ): Mass continuity will read: These equations could be (slightly) compacted by, for example, factoring 1 r 2 {\textstyle {\frac {1}{r^{2}}}} from the viscous terms. However, doing so would undesirably alter the structure of the Laplacian and other quantities.
Navier–Stokes equations use in games:
The Navier–Stokes equations are used extensively in video games in order to model a wide variety of natural phenomena. Simulations of small-scale gaseous fluids, such as fire and smoke, are often based on the seminal paper "Real-Time Fluid Dynamics for Games" by Jos Stam, which elaborates one of the methods proposed in Stam's earlier, more famous paper "Stable Fluids" from 1999. Stam proposes stable fluid simulation using a Navier–Stokes solution method from 1968, coupled with an unconditionally stable semi-Lagrangian advection scheme, as first proposed in 1992.
Navier–Stokes equations use in games:
More recent implementations based upon this work run on the game systems graphics processing unit (GPU) as opposed to the central processing unit (CPU) and achieve a much higher degree of performance.
Many improvements have been proposed to Stam's original work, which suffers inherently from high numerical dissipation in both velocity and mass.
An introduction to interactive fluid simulation can be found in the 2007 ACM SIGGRAPH course, Fluid Simulation for Computer Animation.
General references:
Acheson, D. J. (1990), Elementary Fluid Dynamics, Oxford Applied Mathematics and Computing Science Series, Oxford University Press, ISBN 978-0-19-859679-0 Batchelor, G. K. (1967), An Introduction to Fluid Dynamics, Cambridge University Press, ISBN 978-0-521-66396-0 Currie, I. G. (1974), Fundamental Mechanics of Fluids, McGraw-Hill, ISBN 978-0-07-015000-3 V. Girault and P. A. Raviart. Finite Element Methods for Navier–Stokes Equations: Theory and Algorithms. Springer Series in Computational Mathematics. Springer-Verlag, 1986.
General references:
Landau, L. D.; Lifshitz, E. M. (1987), Fluid mechanics, vol. Course of Theoretical Physics Volume 6 (2nd revised ed.), Pergamon Press, ISBN 978-0-08-033932-0, OCLC 15017127 Polyanin, A. D.; Kutepov, A. M.; Vyazmin, A. V.; Kazenin, D. A. (2002), Hydrodynamics, Mass and Heat Transfer in Chemical Engineering, Taylor & Francis, London, ISBN 978-0-415-27237-7 Rhyming, Inge L. (1991), Dynamique des fluides, Presses polytechniques et universitaires romandes Smits, Alexander J. (2014), A Physical Introduction to Fluid Mechanics, Wiley, ISBN 0-47-1253499 Temam, Roger (1984): Navier–Stokes Equations: Theory and Numerical Analysis, ACM Chelsea Publishing, ISBN 978-0-8218-2737-6 White, Frank M. (2006), Viscous Fluid Flow, McGraw-Hill, ISBN 978-0-07-124493-0 | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Laser turntable**
Laser turntable:
A laser turntable (or optical turntable) is a phonograph that plays standard LP records (and other gramophone records) using laser beams as the pickup instead of using a stylus as in conventional turntables. Although these turntables use laser pickups, the same as Compact Disc players, the signal remains in the analog realm and is never digitized.
History:
William K. Heine presented a paper "A Laser Scanning Phonograph Record Player" to the 57th Audio Engineering Society (AES) convention in May 1977. The paper details a method developed by Heine that employs a single 2.2 mW helium–neon laser for both tracking a record groove and reproducing the stereo audio of a phonograph in real time. In development since 1972, the working prototype was named the "LASERPHONE", and the methods it used for playback was awarded U.S. Patent 3,992,593 on 16 November 1976. Heine concluded in his paper that he hoped his work would increase interest in using lasers for phonographic playback.
History:
Finial Four years later in 1981 Robert S. Reis, a graduate student in engineering at Stanford University, wrote his master's thesis on "An Optical Turntable". In 1983 he and fellow Stanford electrical engineer Robert E. Stoddard founded Finial Technology to develop and market a laser turntable, raising $7 million in venture capital. In 1984 servo-control expert Robert N. Stark joined the effort.A non-functioning mock-up of the proposed Finial turntable was shown at the 1984 Consumer Electronics Show (CES), generating much interest and a fair amount of mystery, since the patents had not yet been granted and the details had to be kept secret. The first working model, the Finial LT-1 (Laser Turntable-1), was completed in time for the 1986 CES. The prototype revealed an interesting flaw of laser turntables: they are so accurate that they "play" every particle of dirt and dust on the record, instead of pushing them aside as a conventional stylus would. The non-contact laser pickup does have the advantages of eliminating record wear, tracking noise, turntable rumble and feedback from the speakers, but the sound is still that of an LP turntable rather than a Compact Disc. The projected $2,500 street price (later raised to $3,786 in 1988) limited the potential market to professionals (libraries, radio stations and archivists) and a few well-heeled audiophiles.The Finial turntable never went into production. After Finial showed a few hand-built (and finicky) prototypes, tooling delays, component unavailability (in the days before cheap lasers), marketing blunders, and high development costs kept pushing back the release date. The long development of the laser turntable exactly coincided with two major events, the early 1980s recession, and the introduction of the Digital Compact Disc, which soon began flooding the market at prices comparable to LPs (with CD players in the $300 range). Vinyl record sales plummeted, and many established turntable manufacturers went out of business as a result.
History:
With over US$20 million in venture capital invested, Finial faced a marketing dilemma: forge ahead with a selling price that would be too high for most consumers, or gamble on going into mass production at a much lower price and hope the market would lower costs. Neither seemed viable in a rapidly-shrinking market.
History:
ELP Finally, in late 1989 after almost seven years of research, Finial's investors cut their losses and liquidated the firm, selling the patents to Japanese turntable maker BSR, which became CTI Japan, which in turn created ELP Japan for continued development of the "super-audiophile" turntable. After eight more years of development the laser turntable was finally put on sale in 1997 – twenty years after the initial proposal – as the ELP LT-1XA Laser Turntable, with a list price of US$20,500 (in 2003 the price was lowered to US$10,500).
History:
The turntable, which uses two lasers to read the groove and three more to position the head, does allow one to vary the depth at which the groove is read, possibly bypassing existing record wear. It will not, however, read clear or colored vinyl records. ELP sells built-to-order laser turntables directly to consumers in two versions (LT-basic, and LT-master), at a reported cost (unpublished) of approximately $16,000 for the basic model.
History:
Optora In May 2018, Almedio of Japan, a computer drive manufacturer, presented the Optora ORP-1 optical (laser) turntable at the HIGH END Munich audio show. Few details were provided by the company because, like the 1984 presentation of the Finial turntable, the Optora was a non-working mockup. Company representatives indicated the turntable would use five lasers and be belt-driven, like the ELP. However, after producing some promotional materials (since deleted), a price was never announced and the Optora has not been put on the market. The company's website devoted to the turntable has since been deleted.
Performance:
In a 2008 review of the model ELP LT-1LRC, Jonathan Valin in The Absolute Sound claimed "If I were to describe its presentation in a few words, they would be "pleasant but dull."" He commended the tonal accuracy of playback, but criticized the lack of dynamic range and bass response (limitations of the vinyl records themselves). Records must be wet-cleaned immediately before playback because, says Valin, "Unlike a relatively massive diamond stylus, which plows through a record’s grooves like the prow of a ship, the ELP’s tiny laser-beam styli have next to no mass [sic] and cannot move dust particles out of their way. Any speck of dirt, however minute, is read by the lasers along with the music." Michael Fremer, writing in Stereophile in 2003 noted, "...consider the LT's many pluses: no rumble or background noise of any kind; no cartridge-induced resonances or frequency-response anomalies; no compromise in channel separation (the ELP guarantees channel separation in excess of what the best cutter heads offer); zero tracking or tracing error; no inner-groove distortion; no skating; no adjustments of VTA or azimuth to worry about; no tangency error (like the cutter head itself, the laser pickup is a linear tracker); no record wear; a claimed frequency response of 10Hz–25kHz; and, because the laser beam is less than a quarter the contact area of the smallest elliptical stylus, it can negotiate sections of the engraved waveform that even the smallest stylus misses." But, he notes, all this comes at a cost: "[T]he LT-2XRC's laser pickup was unable to distinguish groove modulations from dirt. Records that sound dead quiet on a conventional turntable could sound as if I was munching potato chips while listening to the ELP. Bummer. There's a solution, of course: a record-cleaning machine. This can't be considered an "accessory" with the LT: it's mandatory. Even new records fresh out of the jacket can sound crunchy." He concludes, "Ironically, if you listen to the music itself, you won't know you're listening to an LP. It's almost like a reel-to-reel tape. Unfortunately, when there is noise, it will always make you aware that you're listening to an LP. That's the confounding thing about this fabulous contraption."
Optical record scanning:
A similar technology is to scan or photograph the grooves of the record, and then reconstruct the sound from the modulation of the groove revealed by the image. Research groups that developed this technology include: IRENE developed by physicists Carl Haber and Vitaliy Fadeyev of the Lawrence Berkeley National Laboratory. Installed in the Library of Congress late in 2006, IRENE (for Image, Reconstruct, Erase Noise, Etc.) uses a camera rotating around the record and taking detailed photographs of the grooves. Software then uses the digital images to reconstruct the sound. In 2018 the system was used to play, for the first time, the only known recording of Alexander Graham Bell's voice. IRENE often produces a large amount of hiss with the recording, but it is very capable of removing pops and clicks produced by imperfections on the record surface.
Optical record scanning:
SAPHIR system developed at INA in 2002 (patented in France in 2004).
VisualAudio developed by the Swiss National Sound Archives and the School of Engineering and Architecture of Fribourg.
A laser beam reflection method was developed by Japanese scientists in Hokkaido university in 1986 for the purpose of reading audio recordings of Ainu language made on fragile wax cylinders. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Hyperprothrombinemia**
Hyperprothrombinemia:
Hyperprothrombinemia is a state of high of prothrombin levels in the blood which leads to hypercoagulability. An example of a genetic cause includes the mutation prothrombin G20210A. Hyperprothrombinemia is a risk factor for venous thromboembolism. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Truncated 120-cells**
Truncated 120-cells:
In geometry, a truncated 120-cell is a uniform 4-polytope formed as the truncation of the regular 120-cell. There are three truncations, including a bitruncation, and a tritruncation, which creates the truncated 600-cell.
Truncated 120-cell:
The truncated 120-cell or truncated hecatonicosachoron is a uniform 4-polytope, constructed by a uniform truncation of the regular 120-cell 4-polytope. It is made of 120 truncated dodecahedral and 600 tetrahedral cells. It has 3120 faces: 2400 being triangles and 720 being decagons. There are 4800 edges of two types: 3600 shared by three truncated dodecahedra and 1200 are shared by two truncated dodecahedra and one tetrahedron. Each vertex has 3 truncated dodecahedra and one tetrahedron around it. Its vertex figure is an equilateral triangular pyramid.
Truncated 120-cell:
Alternate names Truncated 120-cell (Norman W. Johnson) Tuncated hecatonicosachoron / Truncated dodecacontachoron / Truncated polydodecahedron Truncated-icosahedral hexacosihecatonicosachoron (Acronym thi) (George Olshevsky, and Jonathan Bowers) Images
Bitruncated 120-cell:
The bitruncated 120-cell or hexacosihecatonicosachoron is a uniform 4-polytope. It has 720 cells: 120 truncated icosahedra, and 600 truncated tetrahedra. Its vertex figure is a digonal disphenoid, with two truncated icosahedra and two truncated tetrahedra around it.
Alternate names Bitruncated 120-cell / Bitruncated 600-cell (Norman W. Johnson) Bitruncated hecatonicosachoron / Bitruncated hexacosichoron / Bitruncated polydodecahedron / Bitruncated polytetrahedron Truncated-icosahedral hexacosihecatonicosachoron (Acronym Xhi) (George Olshevsky, and Jonathan Bowers) Images
Truncated 600-cell:
The truncated 600-cell or truncated hexacosichoron is a uniform 4-polytope. It is derived from the 600-cell by truncation. It has 720 cells: 120 icosahedra and 600 truncated tetrahedra. Its vertex figure is a pentagonal pyramid, with one icosahedron on the base, and 5 truncated tetrahedra around the sides.
Truncated 600-cell:
Alternate names Truncated 600-cell (Norman W. Johnson) Truncated hexacosichoron (Acronym tex) (George Olshevsky, and Jonathan Bowers) Truncated tetraplex (Conway) Structure The truncated 600-cell consists of 600 truncated tetrahedra and 120 icosahedra. The truncated tetrahedral cells are joined to each other via their hexagonal faces, and to the icosahedral cells via their triangular faces. Each icosahedron is surrounded by 20 truncated tetrahedra. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Oxamate**
Oxamate:
Oxamate is the carboxylate anion of oxamic acid. Oxamate has a molecular formula of C2H2NO3− and is an isosteric form of pyruvate. Salts and esters of oxamic acid are known collectively as oxamates.
Oxamate:
Oxamate is a competitive inhibitor of the enzyme lactate dehydrogenase. Oxamate is a possible pyruvate analog that has the ability to halt lactate production by inhibiting lactate dehydrogenase, effectively stopping the conversation process of pyruvate to lactate.Oxamate, as a lactate dehydrogenase (LDH) inhibitor, plus phenformin, an anti-diabetic agent, has been tested in conjunction with one another and it was shown that this combination has potential anti-cancer properties. Phenformin when administered by itself has a high incidence of lactic acidosis. Due to the inherent ability of oxamate to prevent the conversion of pyruvate to lactate, oxamate can be used to counterbalance the side effects of phenformin.Oxamate also plays inhibiting roles with oxaloacetate, an important intermediate for the citric acid cycle. Oxamate competes and binds to the carboxyl transferase domain active site, and reverses the reaction of oxalaoacetate decarboxylation by pyruvate carboxylase. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Communications training**
Communications training:
Communications training or communication skills training refers to various types of training to develop necessary skills for communication. Effective communication is vital for the success in various situations. Individuals undergo communications training to develop and improve communication skills related to various roles in organizations.
Purpose:
In organizations, it is necessary to communicate with different sub-groups and overcome difficulties in effective communication. Since each sub-group has a unique sub-culture, an effective communications trainer may assist organizational members in improving communications between sub-groups of the organization. It is necessary to ensure that communications between individuals the various sub-cultures serve to meet the mission and goals of the organization. Communications training can assist leaders to develop the ability to perceive how various individuals and subgroups relate to each other and make appropriate interventions
Types of skill development:
Listening skills Influence Skills Responding to conflict Customer service Assertiveness skills Negotiation Facilitation Report writing; business and technical writing Public speaking, effective presentation Speaking skills Interacting skills
Benefits:
Business communication training: It is possible for developing the skills needed for business networking and enhance their communication skills. It helps in communicating the apt message to the appropriate person at the most right time and to effectively manage and develop assertive skills. It enable candidates to manage competently, maintain long-term relationships, form new alliances, meet new people and establish contact with them and develop relationship with them Corporate communications training: It is useful for corporate events and help in dealing with other corporate participants, besides being helpful for routine dealings.
Benefits:
Executive communication training: It focuses on how to conduct meetings by helping to develop facilitation skills and through exceptional executive communication coaching, candidates learn how to open, manage, as well as end meetings.
Crisis communication training: It enables candidates to communicate while dealing with the various difficulties and emergencies that can arise including conflict management and change management. With training, candidates will be fit to come up with beneficial solutions for solving the crisis or conflict or make change/transition easier.
Public speaking training: It is very useful to make presentations, for developing their verbal communication skills so that it is possible to express their facts publicly with great confidence. This is useful for even sales and marketing personnel who need to express things in the best possible way.
Effective Training:
In order to maximize the benefits of instruction, some key points such as management training, identifying your audience, and up to date use of technology can be used to fully profit the managers as well as the members of the organization.
Training for management must be done on a regular basis gives an advantage to any institution since they can provide ongoing feedback to personnel in order to ensure the good function of the different components of an association. Leadership instruction as well as communication skills education are some examples of management training.
Identifying your audience, in this case, the format of the organization such as family business, small business, event, charity group, or simply meetings enables you to apply the required techniques to get the most out of your training and preparation sessions.
Effective Training:
As technology grows, its important to keep your preparation up-to-date by using all means necessary. The Internet, computers as well as E-learning provide new insights to effective training and can be adapted to fit different needs for different companies. It's also very important to get constant feedback from the members as well as having assessment strategies to ensure that the training that is being provided is useful and productive to not waste time and resources.
Effective Training:
In the medical field, recent research draws on available evidence from general educational literature, as well as specific literature on communication skills training (CST). These studies "delineate how educational interventions should be organized in order to enhance clinicians’ communication skills learning and practice. CST interventions need to be learner- and practice-centered and include core conceptual knowledge and experiential opportunities for practice, reflection, feedback, and rehearsal". | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Finger-four**
Finger-four:
The finger-four formation (also known as the "four finger formation" and the "Fingertip Formation") is a flight formation used by fighter aircraft. It consists of four aircraft, and four of these formations can be combined into a squadron formation.
Description:
The formation consists of a flight of four aircraft, consisting of a "lead element" and a "second element", each of two aircraft. When viewing the formation from above, the positions of the planes resemble the tips of the four fingers of a human right hand (without the thumb), giving the formation its name.
Description:
The lead element is made up of the flight leader at the very front of the formation and one wingman to his rear left. The second element is made up of an additional two planes, the element leader and his wingman. The element leader is to the right and rear of the flight leader, followed by the element wingman to his right and rear.
Description:
Both the flight leader and element leader have offensive roles, in that they are the ones to open fire on enemy aircraft while the flight remains intact. Their wingmen have a defensive role — the flight wingman covers the rear of the second element and the element wingman covers the rear of the lead element.
Four of these flights can be assembled to form a squadron formation which consists of two staggered lines of fighters, one in front of the other. Each flight is usually designated by a color (i.e. Red, Blue, Yellow, and Green).
One disadvantage of the finger-four formation was that it left the least-experienced flier, the number two wingman, in the most exposed position (as "tail-end Charlie"). These were particularly vulnerable to a surprise diving attack, and in some cases could be picked off without the others in the formation even noticing.
History:
The formation was developed by several air forces independently in the 1930s. The Finnish Air Force adopted it in 1934 and 1935. Luftwaffe pilots developed the formation independently in 1938 during the Spanish Civil War, and were the first to use it in combat.
History:
Finnish experiment During the 1930s, the Finnish Air Force, aware of its weakness in numbers compared to its neighbours, sought to offset the disadvantage by a radical rethink of its tactics. The new tactical philosophy emphasized aggression, a willingness to attack regardless of the odds, as well as shooting accuracy at a time when aerobatic skill was prized by most air forces. Hand in hand with the changes was the adoption of the pair and finger-four formation, which allowed economy and gave the flexibility that the new tactics required. The aircraft in the new formations had greater vertical and horizontal separation and so they were free to scan in all directions for enemy aircraft, rather than focusing on maintaining a close formation. That allowed the pilots to maintain greater situational awareness and reduced their chance of being spotted by the enemy. The two pairs could split and attack on their own. The pilot who spotted the enemy would become the leader of the pair or even the whole flight for the duration of the attack, as he had the best situational awareness at that moment.
History:
With no guarantee of success, the FAF adopted the new tactics and was later to find the validity of this approach during the Winter War (1939–1940) with the Soviet Union. The Finnish Air Force proved their effectiveness by achieving a 16:1 kill ratio with the finger-four against the Soviet Air Force, which used the conventional Vic formation.
History:
German experience in Spain Involvement in the Spanish Civil War gave the fledgling Luftwaffe an early experience of combat conditions, but the "Condor Legion" quickly found its main fighter aircraft, the He 51, was outclassed by the Soviet I-16 in service with the Republicans. The new Bf 109B was effective but in short supply. For some time, only six were available, making the 3-plane Kette impractical. Flying in pairs (Rotte) or a two-pair Schwarm and using a more open formation (made possible by radio communication between aircraft) was found to confer other benefits.
History:
Most notable in its development and use in the Luftwaffe were Günther Lützow and Werner Mölders and their fellow airmen. In the Luftwaffe, the flight (German: Schwarm) was made up of two pairs (German: Rotte) of aircraft. Each Rotte was composed of a leader and a wingman. The Rottenführer (pair leader) would attack enemy aircraft, leave his wingman to scan for threats and protect him while he engaged the enemy. The Germans thus eschewed the Finnish Air Force's more flexible approach.
History:
The Luftwaffe continued the use of the formation during its Battle of France, the Battle of Britain and the invasion of the Soviet Union in which its effectiveness was shown to be considerably greater than the standard three-aircraft "Vic" close formation used by its opponents. Later, the RAF, the United States Army Air Forces (USAAF) and Soviet Air Forces adapted their tactics.
History:
Other operators The Soviet Air Force units in the Spanish Civil War soon adopted the formation flying against the Germans and in 1938 recommended its use when they returned home. However, most of the Spanish veterans were swept away during Stalin's Great Purge of the armed forces, and the more conservative "Vic" remained the standard Soviet formation. The pary and zveno were not reintroduced until the after Operation Barbarossa forced reforms by Alexander Novikov in 1942 and 1943.The RAF similarly could not radically reform their fighter tactics until the end of the Battle of Britain. The easing of the pressure and a switch to a more offensive stance led to various experiments with formations. The flying ace Douglas Bader was the first RAF pilot to try the formation, in May 1941. After some refining it became the standard formation of his Duxford Wing and eventually spread throughout RAF.The United States Army Air Corps and Naval Aviation began using a concept called "Fighting Pair" from 1940 to 1941. Japan also adopted the finger-four formation during World War II.
Missing man formation:
The finger-four formation became less common after World War II. However, it is still used in the "Missing Man Formation" at pilots' funeral ceremonies. The formation performs a fly-by in level flight over the funeral, at which point the second element leader climbs vertically and departs the formation, symbolizing the departure of the person being honored. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Nipa palm vinegar**
Nipa palm vinegar:
Nipa palm vinegar, also known as sukang sasâ or sukang nipa, is a traditional Filipino vinegar made from the sap of the nipa palm (Nypa fruticans). It is one of the four main types of vinegars in the Philippines, along with coconut vinegar, cane vinegar, and kaong palm vinegar. It is usually sold under the generic label of "palm vinegar".Nipa palm vinegar is listed in the Ark of Taste international catalogue of endangered heritage foods by the Slow Food movement. Along with other traditional vinegars in the Philippines, it is threatened by the increasing use of industrially-produced vinegars.
Names:
Nipa palm vinegar is known as sukang sasa or sukang nipa in native languages in the Philippines. Both nipa and sasa are the native names of the nipa palm in Tagalog; while sukâ (with the Tagalog enclitic suffix -ng) means "vinegar". It is also known as sukang Paombong after the town of Paombong, Bulacan where it is a traditional industry. The name of the town itself is allegedly from Tagalog bumbóng ("bamboo tube"), the main equipment in gathering nipa sap before plastic or glass containers became prevalent. It is also sometimes known as sukang tubâ, from tubâ, the general term for palm toddy produced from various palm trees in the Philippines, including coconut, buri palm (Corypha elata), and kaong palm (Arenga pinnata).
Traditional production:
Nipa palm vinegar is gathered from mature nipa palms that grow in muddy soil beside brackish rivers and estuaries. The stalk of the nipa palm is cut and a container (traditionally bumbóng, bamboo tubes) is placed underneath to collect sap. The harvesters traditionally shake or kick the base of the leaves as they collect the containers to induce the sap the flow. They may also sometimes bend the stalk. They are collected twice a day as they fill up, though it may take longer during dry seasons.The collected sap are placed in tapayan, large earthen jars traditionally used for fermentation. The sap relies on wild yeast to turn the sugars into ethanol. This turns the sap into a traditional palm toddy called tubâ. Leaving it to ferment further, however, allows Acetobacter from the air to oxidise the ethanol into acetic acid. It is harvested once the level of acidity reaches four or five percent. The length of time it takes to produce nipa palm vinegar ranges from two to three weeks, though it is faster if a starter culture of yeast is used.Nipa palm sap has a relatively high sugar content, containing 15 to 22% sugar. This makes nipa palm vinegar slightly sweeter and less sharp than coconut vinegar. It is also slightly salty due to sodium content of the sap from the habitat of nipa palms. The vinegar when newly made is typically cloudy white. Due to the high iron content of the sap, the vinegar tends to turn orange to dark red as it ages. The vinegar also contains calcium, magnesium, and potassium. The sourness of the vinegar depends on how long it has been allowed to ferment.
Modern production:
The production of nipa palm vinegar is usually associated with the town of Paombong in the province of Bulacan, where it is a prevalent local industry. However, it is also produced in other parts of the Philippines. The production of nipa palm vinegar is labor-intensive and it is predominantly only sold in local markets. Usually in roadside stands or by hawkers along with tubâ palm wines. Along with other traditional vinegars in the Philippines, which also have problems penetrating the national market, it is threatened by the increasing use of industrially-produced vinegars. Many nipa farmers are converting their nipa plantations into fish farms. It is listed in the Ark of Taste international catalogue of endangered heritage foods by the Slow Food movement.
Culinary uses:
Vinegar is one of the most important ingredients in traditional Filipino cuisine. Like other types of vinegars, nipa palm vinegar is used primarily in dipping sauces (sawsawan). It may also be sold spiced with ginger, garlic, and chili peppers (which are boiled beforehand). It can also be used in salad dressings as well as an ingredient in various dishes like paksiw and atchara pickles. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Paterson's worms**
Paterson's worms:
Paterson's worms are a family of cellular automata devised in 1971 by Mike Paterson and John Horton Conway to model the behaviour and feeding patterns of certain prehistoric worms. In the model, a worm moves between points on a triangular grid along line segments, representing food. Its turnings are determined by the configuration of eaten and uneaten line segments adjacent to the point at which the worm currently is. Despite being governed by simple rules the behaviour of the worms can be extremely complex, and the ultimate fate of one variant is still unknown.
Paterson's worms:
The worms were studied in the early 1970s by Paterson, Conway and Michael Beeler, described by Beeler in June 1973, and presented in November 1973 in Martin Gardner's "Mathematical Games" column in Scientific American.Electronic Arts' 1983 game Worms? is an interactive implementation of Paterson's worms, where each time a worm has to turn in a way that it lacks a rule for, it stops and lets the user choose a direction, which sets that rule for that worm.
History:
Paterson's worms are an attempt to simulate the behaviour of prehistoric worms. These creatures fed upon sediment at the bottom of ponds and avoided retracing paths they had already travelled because food would be scarce there but, because food occurred in patches, it was in the worm's interest to stay near previous trails. Different species of worm had different innate rules regarding how close to travelled paths to stay, when to turn, and how sharp a turn to make. In 1969 Raup and Seilacher created computer simulations of the fossilized worm trails, and these simulations inspired Paterson and Conway to develop a simple set of rules to study idealized worms on regular grids.Conway's original model was a worm on an orthogonal grid but this produced only three different species of worm, all with rather uninteresting behaviour. Paterson considered worms on a triangular grid. Paterson's worms were described by Beeler in a Massachusetts Institute of Technology AI Memo (#[1]) and were presented in November 1973 in Martin Gardner's "Mathematical Games" column in Scientific American, and later reprinted in Gardner 1986. These simulations differed in approach from other cellular automata developed around the same time, which focused on cells and the relationships between them. Simple computer models such as these are too abstract to accurately describe the behaviour of the real creatures, but they do demonstrate that even very simple rules can give rise to patterns resembling their tracks.
Rules:
The worm starts at some point of an infinite triangular grid. It starts moving along one of the six gridlines that meet at each point and, once it has travelled one unit of distance, it arrives at a new point. The worm then decides, based on the distribution of traversed and untraversed gridlines, what direction it will take. The directions are relative to the worm's point of view. If the worm has not encountered this exact distribution before it may leave along any untraversed gridline. From then on, if it encounters that distribution again, it must move in the same way. If there are no untraversed gridlines available, the worm dies and the simulation ends.
Discussion:
There are many different types of worm depending on which direction they turn when encountering a new type of intersection. The different varieties of worm can be classified systematically by assigning every direction a number and listing the choice made every time a new type of intersection is encountered.The six directions are numbered as follows: So direction 0 indicates the worm continues to travel straight ahead, direction 1 indicates the worm will make a right turn of 60° and similarly for the other directions. The worm cannot travel in direction 3 because that is the gridline it has just traversed. Thus a worm with rule {1,0,5,1} decides to travel in direction 1 the first time it has to make a choice, in direction 0 the next time it has to make a choice and so on. If there is only one available gridline, the worm has no choice but to take it and this is usually not explicitly listed.
Discussion:
A worm whose ruleset begins with 0 continues in a straight line forever. This is a trivial case, so it is usually stipulated that the worm must turn when it encounters a point with only uneaten gridlines. Furthermore, to avoid mirror-image symmetrical duplicates, the worm's first turn must be a right hand turn. A worm dies if it returns to its origin a third time, because there are then no untraversed edges available. Only the origin can be lethal to the worm.There are 1,296 possible combinations of worm rules. This can be seen by the following argument: If the worm encounters a node with no eaten segments, other than the one it has just eaten, it can either make a sharp turn or a gentle one. This is the situation shown in the figure above. Since the initial choice of left or right produce combinations that simply mirrors of each other, they are not effectively different.
Discussion:
If it encounters a node with one eaten segment, it can leave along any of the remaining four. Only the worm's first return to the origin has this character.
For two eaten segments, the location of the eaten segments is important. The only type of two-segment intersections that can exist is that produced by the first rule, for which there are four distinct approach directions, each of which offers a choice of three departure directions. This allows for 81 different alternatives in choosing rules.
If the worm returns to the origin, it will encounter three eaten segments and must choose between the two remaining uneaten ones regardless of their distribution.
Discussion:
For four eaten segments, there is only one uneaten segment left and the worm must take it.There are therefore 2×4×81×2x1=1,296 different combinations of rules. Many of these are mirror-image duplicates of others, and others die before having to make all the choices in their ruleset, leaving 411 distinct species (412 if the infinite straight-line worm is included). 336 of these species eventually die. 73 patterns exhibit infinite behaviour, that is, they settle into a repeating pattern that does not return to the origin. A further two are strongly believed to be infinite and one remains unsolved. Eleven of the rules exhibit complicated behaviour. They do not die even after many billions of iterations, nor do they adopt an obviously infinite pattern. Their ultimate fate was unknown until 2003 when Benjamin Chaffin developed new methods of solving them. After many hours of computer time, nine of the eleven rules were solved, leaving the worms with rules {1,0,4,2,0,2,0} and {1,0,4,2,0,1,5}. The first of these was solved by Tomas Rokicki, who determined that it halts after 57 trillion (5.7×1013) timesteps, leaving only {1,0,4,2,0,1,5} unsolved. According to Rokicki, the worm is still active after 5.2×1019 timesteps. He used an algorithm based on Bill Gosper's Hashlife to simulate the worms at extraordinary speeds. This behaviour is considerably more complex than the related rectangular grid worm, which has a longest path of only 16 segments.It is possible for two different species of worm to produce the same path, though they do not necessarily traverse it in the same order. The most common path is also the shortest: the seven point MOT test/fallout shelter symbol. One example of this path is shown in the animated figure above. In total there are 299 different paths, and 209 of these are produced by just one species. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Prison rape**
Prison rape:
Prison rape or jail rape is sexual assault of people while they are incarcerated. The phrase is commonly used to describe rape of inmates by other inmates, or to describe rape of inmates by staff.
China:
In February 2021, BBC News reported eyewitness accounts of systematic rape of Uyghur women in the Xinjiang internment camps.Multiple women who were formerly detained in the Xinjiang internment camps have publicly made accusations of systemic sexual abuse, including rape. Sayragul Sauytbay, a teacher who was forced to work in the camps, told the BBC that employees of the camp in which she was detained conducted rapes en masse, saying that camp guards "picked the girls and young women they wanted and took them away". She also told the BBC of an organized gang rape, in which a woman around age 21 was forced to make a confession in front of a crowd of 100 other women detained in the camps, before being raped by multiple policemen in front of the assembled crowd. Tursunay Ziawudun, a woman who was detained in the camps for a period of nine months, told the BBC that women were removed from their cells "every night" to be raped by Chinese men, and that she was subjected to three separate instances of gang rape while detained. Qelbinur Sedik, an Uzbek woman from Xinjiang, has stated that Chinese police sexually abused detainees during electric shock tortures, saying that "there were four kinds of electric shock... the chair, the glove, the helmet, and anal rape with a stick".
Iran:
Sexual violence against political prisoners is prevalent in Iran. It is allegedly ignored or even facilitated by authorities.Reports issued to the United Nations allege that rape has been used by interrogators in Iran for decades. During the 1980s, following the Iranian Islamic Revolution, the rape of female political prisoners was so prevalent that it prompted Hussein-Ali Montazeri, Supreme Leader Ayatollah Khomeini's then-deputy, to write the following to Khomeini in a letter dated 7 October 1986: "Did you know that young women are raped in some of the prisons of the Islamic Republic?" Two prominent members of Iran's human rights community, the feminist lawyer and journalist Shadi Sadr and the blogger and activist Mojtaba Saminejad published essays online from inside Iran saying prison rape has a long history in the Islamic Republic.In the 2009 Iranian presidential election protests, opposition groups reported thousands were arrested and tortured in prisons around the country, with former inmates alleging mass rape of men, women and children by the Islamic Revolutionary Guards, in prisons such as Kahrizak and Evin.Following the 2009 presidential election, Iranian presidential candidate Mehdi Karroubi said several protesters held behind bars in Evin Prison had been savagely raped, according to a confidential letter to former president and cleric Akbar Hashemi Rafsanjani. Karroubi said this was a "fragment" of the evidence he had and that if the denials did not stop, he would release even more.On 9 August 2009, in a letter to the Chairman of the Expediency Discernment Council of Iran, Mehdi Karroubi demanded investigation of Iranian prisons for possible torture and, in particular, sexual harassment of men and women. On 19 August, he wrote to parliament speaker Ali Larijani, asking to meet with him, President Mahmoud Ahmadinejad, judiciary chief Ayatollah Sadeq Larijani, former president Akbar Hashemi Rafsanjani and the state prosecutor to "personally present my documents and evidence over the cases of sexual abuse in some prisons specially Kahrizak." Ali Larijani and Sadeq Larijani (Judiciary committee) both officially rejected his claims and Ali Khamenei's representatives, and Vice Chairman of National Security Commission of the parliament demanded Karroubi's arrest.
Turkey:
Human Rights Watch and Amnesty International have both released reports of widespread rape and abuse of prisoners in Turkey spanning multiple decades. Kurdish prisoners have also been specifically targeted for rape and other forms of sexual violence.
Middle East:
Rape is regularly used in prisons across the wider Middle East. Sexual abuse of detained women, children and men is rampant in UAE, Saudi and Bahraini prisons.
United States:
In the United States, the overwhelming majority of prison rape cases involve men who are raped by other men. This is due in part to the fact that in the United States the vast majority of incarcerated people are men. Sexual contact with inmates by prison staff is illegal, regardless of supposed consent.Public awareness of common prison rape is a relatively recent development, and estimates of its prevalence have varied widely over the past several decades. In 1974, Carl Weiss and David James Friar wrote that 46 million Americans would one day be incarcerated; of that number, they held that 10 million would be raped.According to a US Department of Justice report from 2013, an estimated 5.0% of people incarcerated in state and federal prison, and 3.2% of those in jail, reported at least one incident of sexual victimization in the prior 12 months. However, advocates dispute the accuracy of the numbers due to under-reporting of sexual assaults in prison, especially among incarcerated youths.In terms of individuals' risk over their entire incarceration, estimates from the 1980s and 1990s range widely. A 1992 estimate from the Federal Bureau of Prisons suggested that between 9% and 20% of inmates had been sexually assaulted. Similarly, studies from 1982 and 1996, concluded that the rate was somewhere between 12% and 14%. In New York State maximum security prisons, a 1986 study put the proportion at around 23%. By contrast, Christine Saum's 1994 survey of 101 inmates determined that 5 had been sexually assaulted.The Prison Rape Elimination Act of 2003 was the first United States federal law passed specifically dealing with the sexual assault of prisoners. The bill was signed into law on 4 September 2003. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**1,2-Diformylhydrazine**
1,2-Diformylhydrazine:
1,2-Diformylhydrazine is the chemical compound with the formula N2H2(CHO)2. It is a white, water-soluble solid. A related species is the monoformyl analog, called formic hydrazide (HCON2H3, 624-84-0 ).
As verified by X-ray crystallography, it is a planar molecule with N-N, N-C, and C=O distances of 1.38, 1.33 and 1.24 Å, respectively. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**1889 Baltimore Orioles season**
Player stats:
Batting Starters by position Note: Pos = Position; G = Games played; AB = At bats; H = Hits; Avg. = Batting average; HR = Home runs; RBI = Runs batted in Other batters Note: G = Games played; AB = At bats; H = Hits; Avg. = Batting average; HR = Home runs; RBI = Runs batted in Pitching Starting pitchers Note: G = Games pitched; IP = Innings pitched; W = Wins; L = Losses; ERA = Earned run average; SO = Strikeouts | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Adjective**
Adjective:
An adjective (abbreviated adj.) is a word that describes a noun or noun phrase. Its semantic role is to change information given by the noun.
Traditionally, adjectives were considered one of the main parts of speech of the English language, although historically they were classed together with nouns. Nowadays, certain words that usually had been classified as adjectives, including the, this, my, etc., typically are classed separately, as determiners.
Here are some examples: That's a funny idea. (attributive) That idea is funny. (predicative) Tell me something funny. (postpositive) The good, the bad, and the funny. (substantive)
Etymology:
Adjective comes from Latin nōmen adjectīvum, a calque of Ancient Greek: ἐπίθετον ὄνομα (surname), romanized: epítheton ónoma, lit. 'additional noun' (whence also English epithet). In the grammatical tradition of Latin and Greek, because adjectives were inflected for gender, number, and case like nouns (a process called declension), they were considered a type of noun. The words that are today typically called nouns were then called substantive nouns (nōmen substantīvum). The terms noun substantive and noun adjective were formerly used in English but are now obsolete.
Types of use:
Depending on the language, an adjective can precede a corresponding noun on a prepositive basis or it can follow a corresponding noun on a postpositive basis. Structural, contextual, and style considerations can impinge on the pre-or post-position of an adjective in a given instance of its occurrence. In English, occurrences of adjectives generally can be classified into one of three categories: Prepositive adjectives, which are also known as "attributive adjectives", occur on an antecedent basis within a noun phrase. For example: "I put my happy kids into the car", wherein happy occurs on an antecedent basis within the my happy kids noun phrase, and therefore functions in a prepositive adjective.
Types of use:
Postpositive adjectives can occur: (a) immediately subsequent to a noun within a noun phrase, e.g. "The only room available cost twice what we expected"; (b) as linked via a copula or other linking mechanism subsequent to a corresponding noun or pronoun; for example: "My kids are happy", wherein happy is a predicate adjective (see also: Predicative expression, Subject complement); or (c) as an appositive adjective within a noun phrase, e.g. "My kids, [who are] happy to go for a drive, are in the back seat." Nominalized adjectives, which function as nouns. One way this happens is by eliding a noun from an adjective-noun noun phrase, whose remnant thus is a nominalization. In the sentence, "I read two books to them; he preferred the sad book, but she preferred the happy", happy is a nominalized adjective, short for "happy one" or "happy book". Another way this happens is in phrases like "out with the old, in with the new", where "the old" means "that which is old" or "all that is old", and similarly with "the new". In such cases, the adjective may function as a mass noun (as in the preceding example). In English, it may also function as a plural count noun denoting a collective group, as in "The meek shall inherit the Earth", where "the meek" means "those who are meek" or "all who are meek".
Distribution:
Adjectives feature as a part of speech (word class) in most languages. In some languages, the words that serve the semantic function of adjectives are categorized together with some other class, such as nouns or verbs. In the phrase "a Ford car", "Ford" is unquestionably a noun but its function is adjectival: to modify "car". In some languages adjectives can function as nouns: for example, the Spanish phrase "un rojo" means "a red [one]".
Distribution:
As for "confusion" with verbs, rather than an adjective meaning "big", a language might have a verb that means "to be big" and could then use an attributive verb construction analogous to "big-being house" to express what in English is called a "big house". Such an analysis is possible for the grammar of Standard Chinese, for example.
Distribution:
Different languages do not use adjectives in exactly the same situations. For example, where English uses "to be hungry" (hungry being an adjective), Dutch, French, and Spanish use "honger hebben", "avoir faim", and "tener hambre" respectively (literally "to have hunger", the words for "hunger" being nouns). Similarly, where Hebrew uses the adjective זקוק (zaqūq, roughly "in need of"), English uses the verb "to need".
Distribution:
In languages that have adjectives as a word class, it is usually an open class; that is, it is relatively common for new adjectives to be formed via such processes as derivation. However, Bantu languages are well known for having only a small closed class of adjectives, and new adjectives are not easily derived. Similarly, native Japanese adjectives (i-adjectives) are considered a closed class (as are native verbs), although nouns (an open class) may be used in the genitive to convey some adjectival meanings, and there is also the separate open class of adjectival nouns (na-adjectives).
Adverbs:
Many languages (including English) distinguish between adjectives, which qualify nouns and pronouns, and adverbs, which mainly modify verbs, adjectives, or other adverbs. Not all languages make this exact distinction; many (including English) have words that can function as either. For example, in English, fast is an adjective in "a fast car" (where it qualifies the noun car) but an adverb in "he drove fast" (where it modifies the verb drove).
Adverbs:
In Dutch and German, adjectives and adverbs are usually identical in form and many grammarians do not make the distinction, but patterns of inflection can suggest a difference: Eine kluge neue Idee.
A clever new idea.
Eine klug ausgereifte Idee.
Adverbs:
A cleverly developed idea.A German word like klug ("clever(ly)") takes endings when used as an attributive adjective but not when used adverbially. Whether these are distinct parts of speech or distinct usages of the same part of speech is a question of analysis. It can be noted that, while German linguistic terminology distinguishes adverbiale from adjektivische Formen, German refers to both as Eigenschaftswörter ("property words").
Determiners:
Linguists today distinguish determiners from adjectives, considering them to be two separate parts of speech (or lexical categories). Determiners formerly were considered to be adjectives in some of their uses. Determiners function neither as nouns nor pronouns but instead characterize a nominal element within a particular context. They generally do this by indicating definiteness (a vs. the), quantity (one vs. some vs. many), or another such property.
Adjective phrases:
An adjective acts as the head of an adjective phrase or adjectival phrase (AP). In the simplest case, an adjective phrase consists solely of the adjective; more complex adjective phrases may contain one or more adverbs modifying the adjective ("very strong"), or one or more complements (such as "worth several dollars", "full of toys", or "eager to please"). In English, attributive adjective phrases that include complements typically follow the noun that they qualify ("an evildoer devoid of redeeming qualities").
Other modifiers of nouns:
In many languages (including English) it is possible for nouns to modify other nouns. Unlike adjectives, nouns acting as modifiers (called attributive nouns or noun adjuncts) usually are not predicative; a beautiful park is beautiful, but a car park is not "car". The modifier often indicates origin ("Virginia reel"), purpose ("work clothes"), semantic patient ("man eater") or semantic subject ("child actor"); however, it may generally indicate almost any semantic relationship. It is also common for adjectives to be derived from nouns, as in boyish, birdlike, behavioral (behavioural), famous, manly, angelic, and so on.
Other modifiers of nouns:
In Australian Aboriginal languages, the distinction between adjectives and nouns is typically thought weak, and many of the languages only use nouns--or nouns with a limited set of adjective-deriving affixes--to modify other nouns. In languages that have a subtle adjective-noun distinction, one way to tell them apart is that a modifying adjective can come to stand in for an entire elided noun phrase, while a modifying noun cannot. For example, in Bardi, the adjective moorrooloo 'little' in the phrase moorrooloo baawa 'little child' can stand on its own to mean 'the little one,' while the attributive noun aamba 'man' in the phrase aamba baawa 'male child' cannot stand for the whole phrase to mean 'the male one.' In other languages, like Warlpiri, nouns and adjectives are lumped together beneath the nominal umbrella because of their shared syntactic distribution as arguments of predicates. The only thing distinguishing them is that some nominals seem to semantically denote entities (typically nouns in English) and some nominals seem to denote attributes (typically adjectives in English).Many languages have participle forms that can act as noun modifiers either alone or as the head of a phrase. Sometimes participles develop into functional usage as adjectives. Examples in English include relieved (the past participle of relieve), used as an adjective in passive voice constructs such as "I am so relieved to see you". Other examples include spoken (the past participle of speak) and going (the present participle of go), which function as attribute adjectives in such phrases as "the spoken word" and "the going rate".
Other modifiers of nouns:
Other constructs that often modify nouns include prepositional phrases (as in "a rebel without a cause"), relative clauses (as in "the man who wasn't there"), and infinitive phrases (as in "a cake to die for"). Some nouns can also take complements such as content clauses (as in "the idea that I would do that"), but these are not commonly considered modifiers. For more information about possible modifiers and dependents of nouns, see Components of noun phrases.
Order:
In many languages, attributive adjectives usually occur in a specific order. In general, the adjective order in English can be summarised as: opinion, size, age or shape, colour, origin, material, purpose. Other language authorities, like the Cambridge Dictionary, state that shape precedes rather than follows age.Determiners and postdeterminers—articles, numerals, and other limiters (e.g. three blind mice)—come before attributive adjectives in English. Although certain combinations of determiners can appear before a noun, they are far more circumscribed than adjectives in their use—typically, only a single determiner would appear before a noun or noun phrase (including any attributive adjectives).
Order:
Opinion – limiter adjectives (e.g. a real hero, a perfect idiot) and adjectives of subjective measure (e.g. beautiful, interesting) or value (e.g. good, bad, costly) Size – adjectives denoting physical size (e.g. tiny, big, extensive) Shape or physical quality – adjectives describing more detailed physical attributes than overall size (e.g. round, sharp, swollen, thin) Age – adjectives denoting age (e.g. young, old, new, ancient, six-year-old) Colour – adjectives denoting colour or pattern (e.g. white, black, pale, spotted) Origin – denominal adjectives denoting source (e.g. Japanese, volcanic, extraterrestrial) Material – denominal adjectives denoting what something is made of (e.g., woollen, metallic, wooden) Qualifier/purpose – final limiter, which sometimes forms part of the (compound) noun (e.g., rocking chair, hunting cabin, passenger car, book cover)This means that, in English, adjectives pertaining to size precede adjectives pertaining to age ("little old", not "old little"), which in turn generally precede adjectives pertaining to colour ("old white", not "white old"). So, one would say "One (quantity) nice (opinion) little (size) old (age) round (shape) [or round old] white (colour) brick (material) house." When several adjectives of the same type are used together, they are ordered from general to specific, like "lovely intelligent person" or "old medieval castle".This order may be more rigid in some languages than others; in some, like Spanish, it may only be a default (unmarked) word order, with other orders being permissible. Other languages, such as Tagalog, follow their adjectival orders as rigidly as English.
Order:
The normal adjectival order of English may be overridden in certain circumstances, especially when one adjective is being fronted. For example, the usual order of adjectives in English would result in the phrase "the bad big wolf" (opinion before size), but instead, the usual phrase is "the big bad wolf".
Order:
Owing partially to borrowings from French, English has some adjectives that follow the noun as postmodifiers, called postpositive adjectives, as in time immemorial and attorney general. Adjectives may even change meaning depending on whether they precede or follow, as in proper: They live in a proper town (a real town, not a village) vs. They live in the town proper (in the town itself, not in the suburbs). All adjectives can follow nouns in certain constructions, such as tell me something new.
Comparison (degrees):
In many languages, some adjectives are comparable and the measure of comparison is called degree. For example, a person may be "polite", but another person may be "more polite", and a third person may be the "most polite" of the three. The word "more" here modifies the adjective "polite" to indicate a comparison is being made, and "most" modifies the adjective to indicate an absolute comparison (a superlative).
Comparison (degrees):
Among languages that allow adjectives to be compared, different means are used to indicate comparison. Some languages do not distinguish between comparative and superlative forms. Other languages allow adjectives to be compared but do not have a special comparative form of the adjective. In such cases, as in some Australian Aboriginal languages, case-marking, such as the ablative case may be used to indicate one entity has more of an adjectival quality than (i.e. from—hence ABL) another. Take the following example in Bardi: In English, many adjectives can be inflected to comparative and superlative forms by taking the suffixes "-er" and "-est" (sometimes requiring additional letters before the suffix; see forms for far below), respectively: "great", "greater", "greatest" "deep", "deeper", "deepest"Some adjectives are irregular in this sense: "good", "better", "best" "bad", "worse", "worst" "many", "more", "most" (sometimes regarded as an adverb or determiner) "little", "less", "least"Some adjectives can have both regular and irregular variations: "old", "older", "oldest" "far", "farther", "farthest"also "old", "elder", "eldest" "far", "further", "furthest"Another way to convey comparison is by incorporating the words "more" and "most". There is no simple rule to decide which means is correct for any given adjective, however. The general tendency is for simpler adjectives and those from Anglo-Saxon to take the suffixes, while longer adjectives and those from French, Latin, or Greek do not—but sometimes sound of the word is the deciding factor.
Comparison (degrees):
Many adjectives do not naturally lend themselves to comparison. For example, some English speakers would argue that it does not make sense to say that one thing is "more ultimate" than another, or that something is "most ultimate", since the word "ultimate" is already absolute in its semantics. Such adjectives are called non-comparable or absolute. Nevertheless, native speakers will frequently play with the raised forms of adjectives of this sort. Although "pregnant" is logically non-comparable (either one is pregnant or not), one may hear a sentence like "She looks more and more pregnant each day". Likewise "extinct" and "equal" appear to be non-comparable, but one might say that a language about which nothing is known is "more extinct" than a well-documented language with surviving literature but no speakers, while George Orwell wrote, "All animals are equal, but some animals are more equal than others". These cases may be viewed as evidence that the base forms of these adjectives are not as absolute in their semantics as is usually thought.
Comparison (degrees):
Comparative and superlative forms are also occasionally used for other purposes than comparison. In English comparatives can be used to suggest that a statement is only tentative or tendential: one might say "John is more the shy-and-retiring type", where the comparative "more" is not really comparing him with other people or with other impressions of him, but rather, could be substituting for "on the whole" or "more so than not". In Italian, superlatives are frequently used to put strong emphasis on an adjective: bellissimo means "most beautiful", but is in fact more commonly heard in the sense "extremely beautiful".
Restrictiveness:
Attributive adjectives and other noun modifiers may be used either restrictively (helping to identify the noun's referent, hence "restricting" its reference) or non-restrictively (helping to describe a noun). For example: He was a lazy sort, who would avoid a difficult task and fill his working hours with easy ones.Here "difficult" is restrictive – it tells which tasks he avoids, distinguishing these from the easy ones: "Only those tasks that are difficult".
Restrictiveness:
She had the job of sorting out the mess left by her predecessor, and she performed this difficult task with great acumen.Here "difficult" is non-restrictive – it is already known which task it was, but the adjective describes it more fully: "The aforementioned task, which (by the way) is difficult" In some languages, such as Spanish, restrictiveness is consistently marked; for example, in Spanish la tarea difícil means "the difficult task" in the sense of "the task that is difficult" (restrictive), whereas la difícil tarea means "the difficult task" in the sense of "the task, which is difficult" (non-restrictive). In English, restrictiveness is not marked on adjectives but is marked on relative clauses (the difference between "the man who recognized me was there" and "the man, who recognized me, was there" being one of restrictiveness).
Agreement:
In some languages, adjectives alter their form to reflect the gender, case and number of the noun that they describe. This is called agreement or concord. Usually it takes the form of inflections at the end of the word, as in Latin: In Celtic languages, however, initial consonant lenition marks the adjective with a feminine singular noun, as in Irish: Here, a distinction may be made between attributive and predicative usage. In English, adjectives never agree, whereas in French, they always agree. In German, they agree only when they are used attributively, and in Hungarian, they agree only when they are used predicatively:
Semantics:
Semanticist Barbara Partee classifies adjectives semantically as intersective, subsective, or nonsubsective, with nonsubsective adjectives being plain nonsubsective or privative.
Semantics:
An adjective is intersective if and only if the extension of its combination with a noun is equal to the intersection of its extension and that of the noun its modifying. For example, the adjective carnivorous is intersective, given the extension of carnivorous mammal is the intersection of the extensions of carnivorous and mammal (i.e., the set of all mammals who are carnivorous).
Semantics:
An adjective is subsective if and only if the extension of its combination with a noun is a subset of the extension of the noun. For example, the extension of skillful surgeon is a subset of the extension of surgeon, but it is not the intersection of that and the extension of skillful, as that would include (for example) incompetent surgeons who are skilled violinists. All intersective adjectives are subsective, but the term 'subsective' is sometimes used to refer to only those subsective adjectives which are not intersective.
Semantics:
An adjective is privative if and only if the extension of its combination with a noun is disjoint from the extension of the noun. For example, fake is privative because a fake cat is not a cat.
A plain nonsubsective adjective is an adjective that is not subsective or privative. For example, the word possible is this kind of adjective, as the extension of possible murderer overlaps with, but is not included in the extension of murderer (as some, but not all, possible murderers are murderers). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Neutral theory of molecular evolution**
Neutral theory of molecular evolution:
The neutral theory of molecular evolution holds that most evolutionary changes occur at the molecular level, and most of the variation within and between species are due to random genetic drift of mutant alleles that are selectively neutral. The theory applies only for evolution at the molecular level, and is compatible with phenotypic evolution being shaped by natural selection as postulated by Charles Darwin. The neutral theory allows for the possibility that most mutations are deleterious, but holds that because these are rapidly removed by natural selection, they do not make significant contributions to variation within and between species at the molecular level. A neutral mutation is one that does not affect an organism's ability to survive and reproduce. The neutral theory assumes that most mutations that are not deleterious are neutral rather than beneficial. Because only a fraction of gametes are sampled in each generation of a species, the neutral theory suggests that a mutant allele can arise within a population and reach fixation by chance, rather than by selective advantage.The theory was introduced by the Japanese biologist Motoo Kimura in 1968, and independently by two American biologists Jack Lester King and Thomas Hughes Jukes in 1969, and described in detail by Kimura in his 1983 monograph The Neutral Theory of Molecular Evolution. The proposal of the neutral theory was followed by an extensive "neutralist–selectionist" controversy over the interpretation of patterns of molecular divergence and gene polymorphism, peaking in the 1970s and 1980s.
Neutral theory of molecular evolution:
Neutral theory is frequently used as the null hypothesis, as opposed to adaptive explanations, for describing the emergence of morphological or genetic features in organisms and populations. This has been suggested in a number of areas, including in explaining genetic variation between populations of one nominal species, the emergence of complex subcellular machinery, and the convergent emergence of several typical microbial morphologies.
Origins:
While some scientists, such as Freese (1962) and Freese and Yoshida (1965), had suggested that neutral mutations were probably widespread, the original mathematical derivation of the theory had been published by R.A. Fisher in 1930. Fisher, however, gave a reasoned argument for believing that, in practice, neutral gene substitutions would be very rare. A coherent theory of neutral evolution was first proposed by Motoo Kimura in 1968 and by King and Jukes independently in 1969. Kimura initially focused on differences among species; King and Jukes focused on differences within species.
Origins:
Many molecular biologists and population geneticists also contributed to the development of the neutral theory. The principles of population genetics, established by J.B.S. Haldane, R.A. Fisher, and Sewall Wright, created a mathematical approach to analyzing gene frequencies that contributed to the development of Kimura's theory.
Haldane's dilemma regarding the cost of selection was used as motivation by Kimura. Haldane estimated that it takes about 300 generations for a beneficial mutation to become fixed in a mammalian lineage, meaning that the number of substitutions (1.5 per year) in the evolution between humans and chimpanzees was too high to be explained by beneficial mutations.
Functional constraint:
The neutral theory holds that as functional constraint diminishes, the probability that a mutation is neutral rises, and so should the rate of sequence divergence.
Functional constraint:
When comparing various proteins, extremely high evolutionary rates were observed in proteins such as fibrinopeptides and the C chain of the proinsulin molecule, which both have little to no functionality compared to their active molecules. Kimura and Ohta also estimated that the alpha and beta chains on the surface of a hemoglobin protein evolve at a rate almost ten times faster than the inside pockets, which would imply that the overall molecular structure of hemoglobin is less significant than the inside where the iron-containing heme groups reside.There is evidence that rates of nucleotide substitution are particularly high in the third position of a codon, where there is little functional constraint. This view is based in part on the degenerate genetic code, in which sequences of three nucleotides (codons) may differ and yet encode the same amino acid (GCC and GCA both encode alanine, for example). Consequently, many potential single-nucleotide changes are in effect "silent" or "unexpressed" (see synonymous or silent substitution). Such changes are presumed to have little or no biological effect.
Quantitative theory:
Kimura also developed the infinite sites model (ISM) to provide insight into evolutionary rates of mutant alleles. If v were to represent the rate of mutation of gametes per generation of N individuals, each with two sets of chromosomes, the total number of new mutants in each generation is 2Nv . Now let k represent the evolution rate in terms of a mutant allele μ becoming fixed in a population.
Quantitative theory:
k=2Nvμ According to ISM, selectively neutral mutations appear at rate μ in each of the 2N copies of a gene, and fix with probability 1/(2N) . Because any of the 2N genes have the ability to become fixed in a population, 1/2N is equal to μ , resulting in the rate of evolutionary rate equation: k=v This means that if all mutations were neutral, the rate at which fixed differences accumulate between divergent populations is predicted to be equal to the per-individual mutation rate, independent of population size. When the proportion of mutations that are neutral is constant, so is the divergence rate between populations. This provides a rationale for the molecular clock, which predated neutral theory. The ISM also demonstrates a constancy that is observed in molecular lineages.
Quantitative theory:
This stochastic process is assumed to obey equations describing random genetic drift by means of accidents of sampling, rather than for example genetic hitchhiking of a neutral allele due to genetic linkage with non-neutral alleles. After appearing by mutation, a neutral allele may become more common within the population via genetic drift. Usually, it will be lost, or in rare cases it may become fixed, meaning that the new allele becomes standard in the population.
Quantitative theory:
According to the neutral theory of molecular evolution, the amount of genetic variation within a species should be proportional to the effective population size.
The "neutralist–selectionist" debate:
A heated debate arose when Kimura's theory was published, largely revolving around the relative percentages of polymorphic and fixed alleles that are "neutral" versus "non-neutral".
The "neutralist–selectionist" debate:
A genetic polymorphism means that different forms of particular genes, and hence of the proteins that they produce, are co-existing within a species. Selectionists claimed that such polymorphisms are maintained by balancing selection, while neutralists view the variation of a protein as a transient phase of molecular evolution. Studies by Richard K. Koehn and W. F. Eanes demonstrated a correlation between polymorphism and molecular weight of their molecular subunits. This is consistent with the neutral theory assumption that larger subunits should have higher rates of neutral mutation. Selectionists, on the other hand, contribute environmental conditions to be the major determinants of polymorphisms rather than structural and functional factors.According to the neutral theory of molecular evolution, the amount of genetic variation within a species should be proportional to the effective population size. Levels of genetic diversity vary much less than census population sizes, giving rise to the "paradox of variation" . While high levels of genetic diversity were one of the original arguments in favor of neutral theory, the paradox of variation has been one of the strongest arguments against neutral theory.
The "neutralist–selectionist" debate:
There are a large number of statistical methods for testing whether neutral theory is a good description of evolution (e.g., McDonald-Kreitman test), and many authors claimed detection of selection. Some researchers have nevertheless argued that the neutral theory still stands, while expanding the definition of neutral theory to include background selection at linked sites.
Nearly neutral theory:
Tomoko Ohta also emphasized the importance of nearly neutral mutations, in particularly slightly deleterious mutations. The population dynamics of nearly neutral mutations are only slightly different from those of neutral mutations unless the absolute magnitude of the selection coefficient is greater than 1/N, where N is the effective population size in respect of selection. The value of N may therefore affect how many mutations can be treated as neutral and how many as deleterious.
Constructive neutral evolution:
The groundworks for the theory of constructive neutral evolution (CNE) was laid by two papers in the 1990s. Constructive neutral evolution is a theory which suggests that complex structures and processes can emerge through neutral transitions. Although a separate theory altogether, the emphasis on neutrality as a process whereby neutral alleles are randomly fixed by genetic drift finds some inspiration from the earlier attempt by the neutral theory to invoke its importance in evolution. Conceptually, there are two components A and B (which may represent two proteins) which interact with each other. A, which performs a function for the system, does not depend on its interaction with B for its functionality, and the interaction itself may have randomly arisen in an individual with the ability to disappear without an effect on the fitness of A. This present yet currently unnecessary interaction is therefore called an "excess capacity" of the system. However, a mutation may occur which compromises the ability of A to perform its function independently. However, the A:B interaction that has already emerged sustains the capacity of A to perform its initial function. Therefore, the emergence of the A:B interaction "presuppresses" the deleterious nature of the mutation, making it a neutral change in the genome that is capable of spreading through the population via random genetic drift. Hence, A has gained a dependency on its interaction with B. In this case, the loss of B or the A:B interaction would have a negative effect on fitness and so purifying selection would eliminate individuals where this occurs. While each of these steps are individually reversible (for example, A may regain the capacity to function independently or the A:B interaction may be lost), a random sequence of mutations tends to further reduce the capacity of A to function independently and a random walk through the dependency space may very well result in a configuration in which a return to functional independence of A is far too unlikely to occur, which makes CNE a one-directional or "ratchet-like" process. CNE, which does not invoke adaptationist mechanisms for the origins of more complex systems (which involve more parts and interactions contributing to the whole), has seen application in the understanding of the evolutionary origins of the spliceosomal eukaryotic complex, RNA editing, additional ribosomal proteins beyond the core, the emergence of long-noncoding RNA from junk DNA, and so forth. In some cases, ancestral sequence reconstruction techniques have afforded the ability for experimental demonstration of some proposed examples of CNE, as in heterooligomeric ring protein complexes in some fungal lineages.CNE has also been put forwards as the null hypothesis for explaining complex structures, and thus adaptationist explanations for the emergence of complexity must be rigorously tested on a case-by-case basis against this null hypothesis prior to acceptance. Grounds for invoking CNE as a null include that it does not presume that changes offered an adaptive benefit to the host or that they were directionally selected for, while maintaining the importance of more rigorous demonstrations of adaptation when invoked so as to avoid the excessive flaws of adaptationism criticized by Gould and Lewontin.
Empirical evidence for the neutral theory:
One of corollaries of the neutral theory is that the efficiency of positive selection is higher in population or species with higher effective population size. This relationship between the effective population size and selection efficiency was evidenced by genomic studies of species including chimpanzee and human and domesticated species. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Abrocitinib**
Abrocitinib:
Abrocitinib, sold under the brand name Cibinqo, is a medication used for the treatment of atopic dermatitis (eczema). It is a Janus kinase inhibitor and it was developed by Pfizer. It is taken by mouth.The most common side effects include nausea (feeling sick), headache, acne, herpes simplex (viral infection of the mouth or the genitals), increased levels of creatine phosphokinase in the blood (an enzyme released into the blood when muscle is damaged), vomiting, dizziness and pain in the upper belly.Abrocitinib was approved for medical use in the European Union in December 2021, and in the United States in January 2022.
Medical uses:
In the EU, abrocitinib is indicated for the treatment of moderate-to-severe atopic dermatitis in adults who are candidates for systemic therapy.In the US, abrocitinib is indicated for the treatment of people twelve years of age and older with refractory, moderate-to-severe atopic dermatitis whose disease is not adequately controlled with other systemic drug products, including biologics, or when use of those therapies is inadvisable.
Medical uses:
Efficacy of abrocitinib in treating atopic dermatitis According to the latest meta-analysis in 2023, abrocitinib is both efficient and safe in treating moderate-to-severe AD in adolescents and adults. It also relieves itching rapidly and alleviates symptoms of AD. Abrocitinib gives significantly better results than the placebo at both 100 mg and 200 mg. The severity of AD is quantified through Eczema Area and Severity Index (EASI), which is based on the severity of lesion clinical signs. Abrocitinib was more effective than placebo in terms of EASI-reduction, but it also decreased other symptoms. The improvement of depression and anxiety was better in the experimental group than that in the control group.
Medical uses:
In another meta-analysis including 2256 patients from three different studies have showed that abrocitinib improved the EASI scores in comparison with dupilumab, even in the second week of treatment. A faster onset of Investigator’s Global Assessment (IGA) response at the second week was also achieved by administrating abrocitinib and early relief of itching occurred at 2 weeks. In other studies, abrocitinib (200 mg dose) achieved rapid relief from itching after four days of treatment compared with dupilumab and placebo in AD patients.
Side effects:
The most common adverse effects in studies were upper respiratory tract infection, headache, nausea, and diarrhea.The most common side effects in clinical trials were nasopharyngitis, nausea, headaches, herpes simplex (including oral herpes, ophthalmic herpes, herpes dermatitis and genital herpes), and increase in blood creatine phosphokinase. Abrocitinib can cause serious infections, malignancy, major cardiac events, thrombosis and other laboratory abnormalities including thrombocytopenia, lymphopenia, and lipid elevations. However, according to clinical data, Abrocitinib is well tolerated. The total adverse reactions weren’t statistically different between the placebo and the dose of 100 mg of abrocitinib. However, it was slightly higher for the dose of 200 mg of abrocitinib. Symptoms such as acne, headache, and nausea, appeared in the first two weeks of starting abrocitinib, and it was not necessary to interrupt the treatment.
Side effects:
In general, the AE frequency of abrocitinib was the same or a little bit higher than in case of placebo or dupilumab.
Pharmacology:
Mechanism of action It is a selective inhibitor of the enzyme janus kinase 1 (JAK1).
Pharmacology:
It inhibits JAK1 by 28 fold of selectivity over JAK2 and more than 340 fold of selectivity over JAK3. Two mechanisms are involved in atopic dermatitis, one involves epidermal barrier disruptions, and the other one is cutaneous inflammation due to the immune system over response. Acute inflammation in AD typically involves IL-13, IL-4, and IL-33. Consequently, inhibiting JAK1 results in suppressing the signaling cytokines IL-4, IL-3, and IL-31. Many other cytokines are involved in AD and mediated by JAK1 such as type II cytokine receptors for IL-22, IL-19, IL-10, IL-20 and glycoprotein 130 (gp130) including IL-6 and IL-12 which are also associated with JAK2 and TYK2; IFN-α and INF-β signal.
Pharmacology:
Pharmacokinetics Abrocitinib is quickly absorbed from the gut and generally reaches highest blood plasma concentrations within one hour. Only 1.0 to 4.4% of the dose are found unmetabolized in the urine. The half-life of abrocitinib is 5 hours and the absorption is not affected by food. A higher dose (400–800 mg) would delay the absorption to 1.5–4 hours. A steady plasma concentration of abrocitinib can be obtained within 48 hours of treatments. The dose is one daily, and abrocitinib is metabolized mainly by cytochrome P450 (CYP450) in liver such as CYP2C9, CYP2C19, CYP3A4 and CYP2B6. The major metabolites of abrocitinib are pyrrolidinone pyrimidine (inactive), 2-hydroxypropyl (active), and 3-hydroxypropyl (active). Dose reduction to half is advisable when abrocitinib is taken with strong inhibitors of CYP2C19. According to phase 1 clinical trials on abrocitinib oral dose of 200 mg, hepatic functions were not altered. However, it is advisable to reduce the dose by half in case of reduced renal function. In serious hepatic impairment and final stages of renal disease, Abrocitinib is contraindicated. Some changes may occur during the abrocitinib treatment such as the reduction in platelet counts after 4 weeks of starting Abrocitinib. However, they will return to normal at the end of the treatment. An increase in LDL, HDL, and total cholesterol levels was also recorded after 4 weeks of Abrocitinib treatment. The increased levels depend on the abrocitinib dose (15% increase in LDL with 200 mg dose versus 10% increase with 100 mg).
History:
April 2016: initiation of Phase 2b trial December 2017: initiation of JADE Mono-1 Phase 3 trial May 2018: Results of Phase 2b trial posted October 2019: Results of Phase 3 trial presented June 2020: Results of second Phase 3 trial publishedThe US Food and Drug Administration (FDA) approved abrocitinib based on evidence from three controlled clinical trials enrolling a total of 1615 participants supporting efficacy and safety. Two of the trials enrolled participants twelve years of age and older with moderate-to-severe atopic dermatitis and one trial enrolled adults with moderate-to-severe atopic dermatitis. The trials were conducted at multiple sites in 18 countries (i.e., United States, Canada, Australia, Mexico, Chile, Great Britain, Poland, Germany, Bulgaria, Hungary, Czech Republic, Latvia, Slovakia, Spain, Italy, Japan, Korea, Taiwan). In addition, safety analyses were performed on the combined results of these 3 controlled clinical trials and one additional controlled study in a total of 1,540 participants.All three trials evaluated two doses of abrocitinib: 100 mg and 200 mg. The monotherapy trials were identically designed, 16-week, randomized, multicenter, double-blind, placebo-controlled, parallel group, phase 3 trials. Trial-AD-3 with concomitant background therapy was a 24-week, multicenter, randomized, double-blind, active-comparator (dupilumab) and placebo-controlled, phase 3 trial. In this trial, at Week 16, subjects previously receiving placebo were re-randomized to receive abrocitinib 100 mg or 200 mg, subjects previously receiving abrocitinib continued on their respective dose and subjects previously receiving dupilumab continued to take placebo.
Society and culture:
Legal status In October 2021, the Committee for Medicinal Products for Human Use (CHMP) of the European Medicines Agency (EMA) adopted a positive opinion, recommending the granting of a marketing authorization for the medicinal product Cibinqo, intended for the treatment of atopic dermatitis. The applicant for this medicinal product is Pfizer Europe MA EEIG. In December 2021, the European Commission approved abrocitinib for the treatment of atopic dermatitis.In January 2022, the US Food and Drug Administration (FDA) approved abrocitinib for adults with moderate-to-severe atopic dermatitis.
Other therapeutic effects of abrocitinib:
In pediatric peanut allergy Abrocitinib has the ability to decrease T-cell activation and the allergen-specific basophil in the case of peanut allergy. Subsequently, the in vitro allergic responses of peanut allergy are reduced.
Abrocitinib may play the role of an immune modulator in oral immunotherapy of peanut or may be administered alone as monotherapy in cases of allergy to certain food.
Other therapeutic effects of abrocitinib:
Therapy of oral lichen planus Oral lichen planus (OLP) is a chronic inflammatory T- cellular disorder that strikes the oral mucosa. In a clinical report in 2022, a fast resolving of OLP was achieved in a patient treated with Abrocitinib. A dose of 200 mg of Abrocitinib was administered daily as monotherapy for twelve weeks. A constant improvement of lesions, a depletion of Wickham striae, and a disappearance of erosions were observed at weeks four and eight of treatment. At week twelve, there was a total recovery of the right buccal mucosa. No adverse events have occurred during the treatment and Abrocitinib was well tolerated by the patient.
Other therapeutic effects of abrocitinib:
Improving the lesions of extensive necrobiosis lipoidica Necrobiosis lipoidica (NL) is chronic granulomatous disease of the skin. It involves shiny patches or plaques with a sclerotic center and inflammatory edge. It may appear on different parts of the body and specially, the front part of the legs. The atrophic scars remain after healing which can be inconvenient for patients. Nevertheless, new lesions may occur. Systemic therapy with abrocitinib was administered at 200 mg/day for 12 weeks and then reduced to 100 mg. A slight stomach ache accompanies the 200 mg dose and no adverse events occurred with the 100 mg dose. An improvement with the old lesions was obvious and no new lesions were observed. The inflammatory edges decreased and the lesions disappeared. Thus, abrocitinib is linked to improving the life quality of the patient. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Second Generation Multiplex**
Second Generation Multiplex:
Second Generation Multiplex is a DNA profiling system used in the United Kingdom to set up the UK National DNA Database in 1995. It is manufactured by ABI (Applied Biosystems).
It contains primers for the following STR (Short Tandem Repeat) loci.
VWA (HUMVWF31/A), D8 (D8S1179), D21 (D21S11), D18 (D18S51), THO (HUMTHO1), FGA (HUMFIBRA) Also contains primers for the Amelogenin sex indicating test.
The primers are tagged with the following fluorescent dyes for detection under electrophoresis.
5-FAM JOE NEDIts use in the United kingdom as the DNA profiling system used by The UK National DNA Database was superseded by the Second Generation Multiplex Plus SGM+ DNA profiling system in 1998 | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Phosphinous acids**
Phosphinous acids:
Phosphinous acids are usually organophosphorus compounds with the formula R2POH. They are pyramidal in structure. Phosphorus is in the oxidation state III. Most phosphinous acids rapidly convert to the corresponding phosphine oxide, which are tetrahedral and are assigned oxidation state V.
Synthesis:
Only one example is known, bis(trifluoromethyl)phosphinous acid, (CF3)2POH. It is prepared in several steps from phosphorus trichloride (Et = ethyl): PCl3 + 2 Et2NH → PCl2NEt2 + Et2NH2Cl 2 P(NEt2)3 + PCl2NEt2 + 2 CF3Br → P(CF3)2NEt2 + 2 BrClP(NEt2)3 P(CF3)2NEt2 + H2O → P(CF3)2OH + HNEt2
Reactions:
With the lone exception of the bis(trifluoromethyl) derivative, the dominant reaction of phosphinous acids is tautomerization: PR2OH → OPR2HEven the pentafluorophenyl compound P(C6F5)2OH is unstable with respect to the phosphine oxide.Although phosphinous acids are rare, their P-bonded coordination complexes are well established, e.g. Mo(CO)5P(OH)3.
Secondary and primary phosphine oxides:
Tertiary phosphine oxides, compounds with the formula R3PO cannot tautomerize. The situation is different for the secondary and primary phosphine oxides, with the respective formulas R2(H)PO and R(H)2PO.
Secondary and primary phosphine oxides | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Alirocumab**
Alirocumab:
Alirocumab, sold under the brand name Praluent, is a medication used as a second-line treatment for high cholesterol for adults whose cholesterol is not controlled by diet and statin treatment. It is a human monoclonal antibody that belongs to a novel class of anti-cholesterol drugs, known as PCSK9 inhibitors, and it was the first such agent to receive FDA approval. The FDA approval was contingent on the completion of further clinical trials to better determine efficacy and safety.Common side effects include nasopharyngitis (cold), injection site reactions, and influenza.It was approved for medical use in the United States and in the European Union in 2015.
Medical uses:
Alirocumab is used as a second-line treatment to lower LDL cholesterol for adults who have a severe form of hereditary high cholesterol and people with atherosclerosis who require additional lowering of LDL cholesterol when diet and statin treatment have not worked. It is administered by subcutaneous injection. As of July 2015, it is not known whether alirocumab prevents early death from cardiovascular disease or prevents heart attacks; a clinical trial to determine outcomes was ongoing at that time, the results of which were expected in 2017. In November 2018, The New England Journal of Medicine published positive results from a clinical trial with alirocumab. According to the study, alirocumab significantly reduced major adverse cardiovascular events by 15% and it was associated with a 15% lower risk of death from any cause (hazard ratio [HR] 0.85; 95% confidence interval [CI], 0.73 to 0.98).In 2021, the U.S. Food and Drug Administration (FDA) added an indication for alirocumab to treat adults with homozygous familial hypercholesterolemia (HoFH), a genetic condition that causes severely high cholesterol. It is not intended to be used alone but instead added to other treatments for HoFH.
Side effects:
Side effects that occurred in more than 2% of people treated with alirocumab in clinical trials and that occurred more frequently than with placebo, included nose and throat irritation, injection site reactions and bruising, flu-like symptoms, urinary tract infection, diarrhea, bronchitis and cough, and muscle pain, soreness, and spasms.There are no available data on use of alirocumab in pregnant women to assess risks to the fetus, nor is there data on use in children.
Pharmacology:
Alirocumab works by inhibiting the PCSK9 protein. PCSK9 binds to the low-density lipoprotein receptor (LDLR) (which takes cholesterol out of circulation), and that binding leads to the receptor being degraded, and less LDL cholesterol being removed from circulation. Inhibiting PCSK9 prevents the receptor from being degraded, and promotes removal of LDL cholesterol from circulation.After subcutaneous administration of alirocumab, maximal suppression of free PCSK9 occurs within 4 to 8 hours and has an apparent half-life of 17 to 20 days. Inhibition is dose-dependent. The antibody is distributed through the circulation, and it is eliminated at low concentrations by binding to its target, and at higher concentrations through a proteolytic pathway.
Chemistry:
Alirocumab is a human monoclonal antibody of the IgG1 isotype. It is made of two disulfide-linked human heavy chains, each disulfide-linked to a human light chain. It has an approximate molecular weight of 146 kDa.It is produced using Chinese hamster ovary cells transfected with recombinant DNA, that are grown in tanks.
History:
The importance of PCSK9 as a biological target for drug discovery emerged in 2003, when a series of discoveries led to identification of the protein and its gene, its role in causing some cases of familial hypercholesterolaemia when some mutations are present, and its role in causing very low levels of LDL cholesterol when other mutations are present.The discovery and validation of the target set off a race among pharmaceutical and biotech companies.Alirocumab was discovered by Regeneron Pharmaceuticals using its "VelocImmune" mouse, in which many of the genes coding for antibodies have been replaced with human genes.: 255–258 In an investor presentation, Regeneron claimed that with their system, it took only about 19 months from when they first immunized mice with PCSK9 until they filed their IND.: Slide 26 Alirocumab was co-developed with Sanofi under a deal made in 2007. Before it received its international nonproprietary name it was known as REGN727 and SAR236553.Phase 1 trial results were reported in 2012 in The New England Journal of Medicine. A phase 3 trial of statin intolerant patients called ODYSSEY ran for 65 weeks. Results were presented at the 2014 European Society of Cardiology meeting. A 78-week study of alirocumab in 2341 people taking statins who were at high risk for cardiovascular events and had high LDL cholesterol levels was published in April 2015. This study showed a significant reduction of LDL cholesterol levels in patients taking both Alirocumab and oral statins compared to placebo patients solely taking oral statins. Studies are ongoing to assess the effects of alirocumab in normocholesterolemic individuals.In July 2014, Regeneron and Sanofi announced that they had purchased a priority review voucher that BioMarin had won for a recent rare disease drug approval for $67.5 million; the voucher cut four months off the regulatory review time for alirocumab and was part of their strategy to beat Amgen to market with the first approval of a PCSK9 inhibitor.In July 2015, the FDA approved alirocumab as a second-line treatment to lower LDL cholesterol for adults who have hereditary high cholesterol and people with atherosclerosis who require additional lowering of LDL cholesterol when diet and statin treatment have not worked. This was the first approval of a PCSK9 inhibitor. The FDA approval was contingent on the completion of further clinical trials to better determine efficacy and safety.Regeneron and Amgen had each filed for patent protection on their monoclonal antibodies and the companies ended up in patent litigation in the U.S. In March 2016, a district court found that alirocumab infringed Amgen's patents; Amgen then requested an injunction barring Regeneron and Sanofi from marketing alirocumab, which was granted in January 2017. The judge gave Regeneron and Sanofi 30 days to appeal before the injunction went into effect. In October, 2017 the US Court of Appeals reversed the ban and ordered a new trial after finding the jury was given improper instructions and evidence was withheld. Regeneron and Sanofi were allowed to continue marketing alirocumab during the appeals process.The U.S. Food and Drug Administration (FDA) granted approval of Praluent to Regeneron Pharmaceuticals, Inc.
Society and culture:
In 2014 as PCSK9 inhibitors approached regulatory approval, market analysts estimated that the overall market for these drugs could be $10B per year, with each of alirocumab and Amgen's competing drugs having sales of $3B per year, and other competitors dividing the remaining $4B, based on estimates of an annual price for alirocumab of $10,000 per year. At the same time, pharmacy benefit managers such as Express Scripts and CVS Caremark, while recognizing that the new drugs could help patients who were otherwise left with uncontrolled cholesterol levels, and recognizing that injectable biopharmaceuticals will always be more expensive than pills, and especially more expensive than generic pills, expressed concerns about the burden of the new costs on the health care system.When the drug was approved in July 2015, the announced price was higher than analysts had predicted: $14,600 a year. Pharmacy benefit managers continued expressing their concerns, as did insurance companies and some doctors, who were especially concerned over the price, in light of the fact that the FDA approval was based on lowering cholesterol alone, and not on better health outcomes, such as fewer heart attacks or longer life.The treatment for people with very high cholesterol that cannot be controlled with diet or statins is apheresis, which is similar to dialysis in that a person visits a clinic each month and his or her blood is mechanically filtered, in this case to remove LDL cholesterol. That treatment costs $8000 per month, or $96,000 per year. The price of alirocumab was determined based in part on making apheresis no longer necessary.
Society and culture:
Names Alirocumab is the International nonproprietary name. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Laplace principle (large deviations theory)**
Laplace principle (large deviations theory):
In mathematics, Laplace's principle is a basic theorem in large deviations theory which is similar to Varadhan's lemma. It gives an asymptotic expression for the Lebesgue integral of exp(−θφ(x)) over a fixed set A as θ becomes large. Such expressions can be used, for example, in statistical mechanics to determining the limiting behaviour of a system as the temperature tends to absolute zero.
Statement of the result:
Let A be a Lebesgue-measurable subset of d-dimensional Euclidean space Rd and let φ : Rd → R be a measurable function with ∫Ae−φ(x)dx<∞.
Then lim log ∫Ae−θφ(x)dx=−essinfx∈Aφ(x), where ess inf denotes the essential infimum. Heuristically, this may be read as saying that for large θ, exp (−θessinfx∈Aφ(x)).
Application:
The Laplace principle can be applied to the family of probability measures Pθ given by Pθ(A)=(∫Ae−θφ(x)dx)/(∫Rde−θφ(y)dy) to give an asymptotic expression for the probability of some event A as θ becomes large. For example, if X is a standard normally distributed random variable on R, then lim log P[εX∈A]=−essinfx∈Ax22 for every measurable set A. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Ethotoin**
Ethotoin:
Ethotoin (previously marketed as Peganone) is an anticonvulsant drug used in the treatment of epilepsy. It is a hydantoin, similar to phenytoin. It is not available in the United States.
Approval history:
1957 Peganone was granted Food and Drug Administration (FDA) approval to Abbott Laboratories for treatment of grand mal (tonic clonic) and partial complex (psychomotor) seizures.
2003 Peganone was acquired from Abbott Laboratories by Ovation Pharmaceuticals (specialty pharmaceutical company who acquire underpromoted branded pharmaceutical products).
2018 It was announced by Recordati Rare Diseases Inc. that due to a combination of low product demand and complex manufacturing difficulties, product manufacturing, distribution and sale was being discontinued.
Indications and usage:
Ethotoin is indicated for tonic-clonic and partial complex seizures.
Dosing:
Ethotoin is available in 250 mg tablets. It is taken orally in 4 to 6 divided doses per day, preferably after food.
Side effects:
Ataxia, visual disturbances, rash and gastrointestinal problems.
Chemistry:
Ethotoin, 3-ethyl-5-phenylimidazolidine-2,4-dione, is synthesized by the reaction of benzaldehyde oxynitrile, with urea or ammonium hydrocarbonate, which forms an intermediate urea derivative which on acidic conditions cyclizes to 5-phenylhydantoin. Alkylation of this product using ethyliodide leads to the formation of ethotoin. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Stratificational linguistics**
Stratificational linguistics:
Stratificational Linguistics, also known as Neurocognitive Linguistics (NCL) or Relational Network Theory (RNT), is an approach to linguistics advocated by Sydney Lamb that suggests language usage and production to be stratificational in nature. It regards the linguistic system of individual speakers as consisting of networks of relations which interconnect across different 'strata' or levels of language. These relational networks are hypothesized to correspond to maps of cortical columns in the human brain. Consequently, Stratificational Linguistics is related to the wider family of cognitive linguistic theories. Furthermore, as a functionalist approach to linguistics, Stratificational Linguistics shares a close relationship with Systemic Functional Linguistics (SFL).
Linguistic strata:
Stratificational Linguistics suggests that the linguistic system may be analyzed according to separate 'strata', or levels. The strata are ordered hierarchically and, whilst there are no clear-cut boundaries between strata, the elements of each stratum share similar characteristics. For example, a lexical item in the lexicogrammatical stratum is typically a specific sequence of phonemes which connects one or more lexical meanings in the semantic stratum. Several strata are involved in the production of a sound from an initial idea. In linguistic production, each stratum provides actualization or realization for the next lower stratum. Thus, speaking a word would involve a realizational pathway from the semantic stratum to the lexicogrammar, then the phonology, and then the phonetics. The reverse direction is true for linguistic perception and comprehension. Some commonly posited stratificational units and their strata include: The phoneme as the unit on the phonemic stratum.
Linguistic strata:
The lexeme as the unit on the lexical or lexicogrammatical stratum.
The morpheme as the unit on the morphemic stratum.
Linguistic strata:
The sememe as the unit on the semantic stratum.In contrast to generativist approaches to linguistics, Stratificational Linguistics does not support the notion of an autonomous stratum for syntax. Instead, the term 'lexicogrammar', borrowed from Systemic Functional Linguistics, is preferred because Stratificational Linguistics suggests that syntactic categories are merely labels for classifying different types of lexemes but do not actually play any role in the realization of the lexemes. Rather, it is posited that what is traditionally called 'syntax' is simply the result of what orderings or sequences of lexemes are possible in the lexicogrammatical system of an individual person. In other words, there is no need to posit a separate stratum for syntax to account for syntactic phenomena. It has been further suggested that each lexeme has its own syntactic pattern which determines how it combines with other lexemes, a stance shared with Construction Grammar.
Relational networks:
Linguistic units in Stratificational Linguistics are conceptualised as relational networks. Simply put, a linguistic unit at any stratum is defined in relation to other units. For example, the phonemic sequence /bɔɪ/ may be analyzed as a network node which is activated when the nodes for /b/, /ɔ/ and /ɪ/ are also activated. Similarly, the node for the sequence /tɔɪ/ gets activated when /t/, /ɔ/ and /ɪ/ are also activated. The two sequences /bɔɪ/ and /tɔɪ/ are defined in relation to the set of phoneme nodes /t/, /b/, /ɔ/ and /ɪ/, and their relationships can be graphed as a relational network diagram. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Ayds**
Ayds:
Ayds Reducing Plan Candy (pronounced as "aids") is a discontinued appetite-suppressant candy that enjoyed strong sales in the 1970s and early 1980s and was originally manufactured by The Carlay Company.
Flavors:
Ayds was available in chocolate, chocolate mint, butterscotch, and caramel flavors, and later a peanut butter flavor was introduced. The original packaging used the phrase "Ayds Reducing Plan vitamin and mineral Candy"; a later version used the phrase "appetite suppressant candy". The active ingredient was originally benzocaine, presumably to reduce the sense of taste to reduce eating, later changed in the candy (as reported by The New York Times) to phenylpropanolamine.
History:
The product was introduced by the Carlay Company of Chicago. A U.S. trademark was registered in 1946 claiming first use in commerce was in 1937.In 1944, the U.S. Federal Trade Commission objected to the claim that the product could cause the user to "lose up to 10 pounds in 5 days, without dieting or exercising".The Carlay Company later became a division of Campana Corporation of Batavia, Illinois. Then Campana bought Allied Laboratories of Kansas City in 1956. Thereafter, Campana was bought by Dow Chemical and its president, Irving Willard Crull, was president of Dow for less than six months, during which time he engineered the sale of Campana to Purex in the 1960s. He again became president of Campana while serving as a vice president of Purex, which allowed Campana to function as a separate division. Crull also relied on socialite and Hollywood friends like Bob Hope and his wife Dolores Hope, Tyrone Power and his wife Linda Christian, and others to promote the Ayds Reducing Plan Candy line. A Cosmopolitan magazine article in November 1956 reported that Crull had already recruited a number of his friends amongst socialite and Hollywood celebrities to promote the Ayds Reducing Plan weight-loss regimen.
History:
In 1981, Purex sold the rights to the Ayds name to Jeffrey Martin Inc. In 1987, Jeffrey Martin, Inc. and its product line (including Ayds Appetite Suppressant and Compoz Sleep Aid) were acquired by the Dep Corporation (sometimes written DEP).
Relations with AIDS:
By the mid-1980s, public awareness of the disease AIDS brought notoriety to the brand due to the phonetic similarity of names and the fact that the disease caused immense weight loss in patients. Initially sales were not negatively affected; in a September 1985 interview, the president of the company that manufactured it stated that, in fact, sales had increased as a result of this connection. Early in 1986, another executive of the manufacturer was quoted, "The product has been around for 45 years. Let the disease change its name."In 1988, by which time the product and its name had been sold to the Dep Corporation, company leadership announced that the company was seeking a new name because sales had dropped as much as 50% due to publicity about the disease. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Google Fusion Tables**
Google Fusion Tables:
Google Fusion Tables was a web service provided by Google for data management. Fusion tables was used for gathering, visualising and sharing data tables. Data are stored in multiple tables that Internet users can view and download.
The web service provided means for visualizing data with pie charts, bar charts, lineplots, scatterplots, timelines, network graphs, HTML-formatted card-layouts, and geographical maps. Data are exported in a comma-separated values file format. Visualizations could be embedded in other websites, and updated realtime as data in the table changed.
From the Fusion Tables website: Google Fusion Tables is a service for data management, integration and collaboration.
You can easily upload data sets from CSV, KML and spreadsheets, and visualize the data using a variety of tools. Users can merge data from multiple tables and conduct detailed discussions about the data (on rows, columns and even cells). You can easily visualize large data sets on Google Maps and embed visualizations on other web pages.
Developers can use our API to build applications over Fusion Tables.
Google closed Fusion Tables on 3 December 2019.
Features:
Fusion Tables accepted a data file structured as a simple database table, typically a .csv but also other delimiters. It also imported KML, reading each KML placemark or geospatial object into its own row. Fusion Tables files were private, unlisted or public, as specified by the user and followed the convention established by other Google Docs apps. Files were then listed and searchable in the user's Google Drive.
Features:
The size of uploaded data set was limited to 250 MB per file with a total limit of 1 GB per user. An API allowed data to be ingested automatically. Visualizations were also embeddable into other web pages to support static or live-updating data within publications.
Features:
Structured Data Search, Publication and Reuse search The 'New file' flow also supported searching on existing published tables, encouraging people to reuse and build on existing data before creating new data or making a new copy of the same data. The 'live update' nature of a re-used table could be an advantage to the user where data sets might be receive corrections or be regularly updated.
Features:
'fusion' The 'fusion' in the name Fusion Tables came from the ability to create a 'file' that is really just a view on a join of two or more other files. For example, to publish a map about election results in Illinois, one could upload a table with election results, and then create another file that joins this table with a KML of US electoral districts. Because it was a virtual join rather than a copy, changes to either of the base tables would be reflected in the joined table. The join would extract the districts relevant to the Illinois elections, and the result would be easy to put on a map and embed in a news article or other website.
Features:
Columns from different tables were displayed with a different background color, to help keep track. Multiple tables could be joined using the same key column. Edits to the data needed to happen to the original underlying table, not in the joined table.
Features:
reuse Fusion Tables encouraged read-only reuse of publicly published data sets, or other data sets shared with the user. Although the user could not edit the read-only data set, the user could create visualizations and filtered views on the data in new tabs in the UI. These views would not affect the original file for the file owner or anyone else, but would appear whenever the user who created them opened the file. These tabs were indicated with a dotted line outline.
Features:
editing The UI supported adding rows and editing data, which was also possible programmatically through the Fusion Tables API.
Data Visualizations During import, Fusion Tables automatically detected various data types in the data, and generated a few appropriate visualizations. All tables saw a row view and a card view; those with many types of location data saw a map visualization automatically created as well.
Data types supported within the table view included standard strings, numbers but also images and KML.
Maps Types of location data automatically detected included: latitude/longitude information in one or two columns, KML place descriptions, and some types of placenames, and addresses, which were sent to the Google Maps Geocoding API in order to put them on the map. The results of geocoding were not available in the table; only on the map.
Features:
Fusion Tables was tightly integrated with the Google Maps geocoding service, as well as the Google Maps API, which supported an experimental FusionTablesLayer. Fusion Tables supported KML descriptions of geographic points, lines and polygons as a datatype within the tables. By providing a way to ingest, manage, merge and style larger quantities of data, Fusion Tables facilitated a blossoming of geographic story-telling. Many data journalists used these features to visualize information acquired through a Freedom of Information Act request as part of their published news stories.
Features:
Card View An HTML subset templating language supported customizable card layout and map infowindows displaying static and data field content. Incorporating a call to the Google Chart API could dynamically render a chart based on the data within a single row in the card or infowindow.
Other visualizations Table view (rows & columns), standard pie charts, scatter plots and line graphs, timeline, chloropleth map, network graph, and motion chart.
Filtering Simple filtering tools provided automatic summaries of values in data columns, and allowed the visualized data to be filtered with checkboxes.
Features:
Publishing and Customizing By supporting simple queries, embeddable HTML snippets for visualizations, and a simple HTML templating language for customizing layouts, Fusion Tables straddled the point-n-click world and the production software engineering world with a 'scriptable' functionality that allowed many data owners with limited software development time or expertise to develop highly custom, expressive websites and tools with their data. See examples.
Features:
Maps created in Fusion Tables could be exported to KML and viewed in Google Earth, making Fusion Tables an important authoring tool for many of the non-profits and NGOs working closely with Google Earth Outreach to spread information about their work.
History & Impact:
Fusion Tables was inspired by the challenges of managing scientific data collections in multi-organization collaborations, such as the DNA barcoding collaboration between University of Pennsylvania ecologist Dan Janzen, the International Barcode of Life, and the University of Guelph.
The website launched as part of Google Labs in June 2009, announced by Alon Halevy and Rebecca Shapley. The service was further described in a scientific paper in 2010.
History & Impact:
Maps Visualization Following positive feedback about Fusion Tables' integration with the Intensity map in the Google Visualization API, the team worked closely with the Google Maps team to add support in Feb 2010 for KML point, line and polygon objects as a native datatype in the tables, visualized on top of Google Maps' basemap. Additionally, some smarts were applied to detect data columns that described locations (like addresses) and to send them to Google's Geocoding service so they could be rendered on the map. Shortly thereafter in May 2010, the FusionTablesLayer was offered as an experimental feature of the Google Maps API.The integration of Fusion Tables with Google Maps through the FusionTablesLayer was Google's first foray into server-side rendering of users' data onto Google Maps tiles. Prior to the FusionTablesLayer, map pins were rendered on top of basemap tiles in the browser client. By creating many objects for the client to track, this could make maps slow, and effectively limited Google Maps to showing approximately 200 user data points. The FusionTablesLayer demonstrated fast, server-side rendering of large and complex user data onto the Google Maps base map.
History & Impact:
The Fusion Tables SQL API supported sending filter queries to the FusionTablesLayer to dynamically adjust the data shown on the map. These maps could be embedded in another webpage with a simple snippet of HTML code. The open-sourced FusionTablesLayer Wizard point-n-click tool helped people create the snippets, and later the snippet was also available easily in the Fusion Tables UI. In May 2011, Fusion Tables added the ability to style (change the color or visual presentation of) data on the map, as well as default and simple HTML customizable infobubbles (shown when an item on the map is clicked) through both the web app and the APIs.Fusion Tables offered a readily accessible solution for working with data on a map that previously required clunky and expensive desktop software. It met many simple GIS use cases. Fusion Tables was presented as part of the Geo track at Google IO in May 2011: Managing and Visualizing your geospatial data with Fusion Tables.
History & Impact:
Adoption In October 2010, FusionTables demonstrated reliability under heavy traffic spikes when hosting the map visualization of the Iraqi War Deaths data set embedded in a news article from The Guardian. Shortly after the March 2011 earthquake and tsunami in Japan, crisis responders used Fusion Tables to reflect road status and shelters with close-to-realtime updates. Google's Crisis Response team continued to use Fusion Tables as a key tool for creating and updating relevant maps after a crisis.In the 2011, as Google Labs was closed, Fusion Tables 'graduated' into the list of default features in Google Docs, under the title "Tables (beta)".In April 2012, Fusion Tables created its own 'labs' track with several experimental features, including a new version of the user interface, a network graph visualization, and a preview of the revised Fusion Tables API, which officially launched in June 2012.Merging tables continued to be a key, if difficult to discover, part of Fusion Tables. Merging tables was, for example, a great way to use publicly available authoritative KML boundaries for places many people might have data about, such as counties or electoral districts. In August 2012, Fusion Tables launched integration with Table Search, another Google Research project from Alon Halevy.
History & Impact:
Presentations & Trainings Fusion Tables was described in talks at the NICAR conference in 2011 and 2013.American Geophysical Union 1 December 2011 - Visualize your data with Google Fusion Tables DigitalNomad - Using Google Fusion Tables and overview deck Reviews Digital Humanities Blog, University of Alabama. Google Fusion Tables.
History & Impact:
Conference Papers & Publications More in Google Scholar Balakrishnan S., Jacqmin-Adams K., Lee H., McChesney R., Shapley R. (2017) Google Fusion Tables. In: Shekhar S., Xiong H., Zhou X. (eds) Encyclopedia of GIS. Springer, Cham Big Data Storytelling through Interactive Maps. Jayant Madhavan, Sreeram Balakrishnan, Kathryn Brisbin, Hector Gonzalez, Nitin Gupta, Alon Halevy, Karen Jacqmin-Adams, Heidi Lam, Anno Langen, Hongrae Lee, Rod McChesney, Rebecca Shapley, Warren Shen. Google Inc. Bulletin of the Technical Committee on Data Engineering, IEEE. June 2012, Vol. 53 No. 2.
History & Impact:
Brambilla M., Ceri S., Cinefra N., Das Sarma A., Forghieri F., Quarteroni S. (2012) Extracting Information from Google Fusion Tables. In: Ceri S., Brambilla M. (eds) Search Computing. Lecture Notes in Computer Science, vol 7538. Springer, Berlin, Heidelberg Field Information Management Systems for DNA Barcoding. Deck John, Gross Joyce, Stones-Havas Steven, Davies Neil, Shapley Rebecca, and Meyer Christopher. Chapter 12 in: W. John Kress and David L. Erickson (eds.), DNA Barcodes: Methods and Protocols, Methods in Molecular Biology, vol. 858 Google Earth and Google Fusion Tables in support of time-critical collaboration: Mapping the deepwater horizon oil spill with the AVIRIS airborne spectrometer. Eliza S. Bradley, Dar A. Roberts, Philip E. Dennison, Robert O. Green, Michael Eastwood, Sarah R. Lundeen, Ian B. McCubbin and Ira Leifer. Earth Science Informatics, 17 September 2011.
History & Impact:
Socializing Data with Google Fusion Tables - Semantic Scholar Fusion Tables: new ways to collaborate on structured data. Jonathan Goldberg Kidon, Massachusetts Institute of Technology. Dept. of Electrical Engineering and Computer Science, 2010. 61 pages.
Google Fusion Tables: Web-Centered Data Management and Collaboration. Hector Gonzalez, Alon Y. Halevy, Christian S. Jensen, Anno Langen, Jayant Madhavan, Rebecca Shapley, Warren Shen, Jonathan Goldberg-Kidon. Proceedings of SIGMOD 2010 (Industrial Track), Indianapolis, Indiana.
Google Fusion Tables: Data Management, Integration and Collaboration in the Cloud. Hector Gonzalez, Alon Y. Halevy, Christian S. Jensen, Anno Langen, Jayant Madhavan, Rebecca Shapley, Warren Shen. Proceedings of the Symposium on Cloud Computing, 2010 (Industrial Track), Indianapolis, Indiana.
Deprecation In December 2018, Google announced that it would retire Fusion Tables on 3 December 2019. An open-source archive tool was created to export existing Fusion Tables maps to an open-sourced visualizer.Fusion Tables had an avid following that was disappointed to learn of the deprecation. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Optical sine theorem**
Optical sine theorem:
In optics, the optical sine theorem states that the products of the index, height, and sine of the slope angle of a ray in object space and its corresponding ray in image space are equal. That is: n0y0a0=niyiai | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Glassy cell carcinoma of the cervix**
Glassy cell carcinoma of the cervix:
Glassy cell carcinoma of the cervix, also glassy cell carcinoma, is a rare aggressive malignant tumour of the uterine cervix. The tumour gets its name from its microscopic appearance; its cytoplasm has a glass-like appearance.
Signs and symptoms:
The signs and symptoms are similar to other cervical cancers and may include post-coital bleeding and/or pain during intercourse (dyspareunia). Early lesions may be completely asymptomatic.
Diagnosis:
The diagnosis is based on tissue examination, e.g. biopsy.Under the microscope, glassy cell carcinoma tumours are composed of cells with a glass-like cytoplasm, typically associated with an inflammatory infiltrate abundant in eosinophils and very mitotically active. PAS staining highlights the plasma membrane.
Treatment:
The treatment is dependent on the stage. Advanced tumours are treated with surgery (radical hysterectomy and bilateral salpingo-opherectomy), radiation therapy and chemotherapy. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**U11 spliceosomal RNA**
U11 spliceosomal RNA:
The U11 snRNA (small nuclear ribonucleic acid) is an important non-coding RNA in the minor spliceosome protein complex, which activates the alternative splicing mechanism. The minor spliceosome is associated with similar protein components as the major spliceosome. It uses U11 snRNA to recognize the 5' splice site (functionally equivalent to U1 snRNA) while U12 snRNA binds to the branchpoint to recognize the 3' splice site (functionally equivalent to U2 snRNA).
Secondary structure:
U11 snRNA has a stem-loop structure with a 5' end as splice site sequence (5' ss) and contains four stem loops structures (I-IV). A structural comparison of U11 snRNA between plants, vertebrates and insects shows that it is folded into a structure with a four-way junction at the 5' site and in a stem loop structure at the 3' site.
Secondary structure:
Binding site during assembly pathway The 5' splice site region possesses sequence complementarity with the 5' splice site of the eukaryotic U12 type pre-mRNA introns. Both the 5' splice site and the Sm binding site are highly conserved in all species. Also, stem loop III is either a possible protein binding site or a base-pairing region since it has a highly conserved nucleotide sequence 'AUCAAGA'.
Role during alternative splicing:
U11 and U12 snRNPs (minor spliceosomal pathway) are functional analogs of U1 and U2 snRNPs (major spliceosomal pathway) whereas the U4 atac/U6 atac snRNPs are similar to U4/U6. Unlike the major splicing pathway, U11 and U12 snRNPs bind to the mRNA as a stable, preformed U11/U12 di-snRNP complex. This is done through the use of seven proteins (65K, 59K, 48K, 35K, 31K, 25K and 20K). Four of them (59K, 48K, 35K and 25K) are associated with U11 snRNA.During the formation of the spliceosome, the 5' end of U11 and U12 snRNAs interact with the 5' splice site and branchpoint sequence of the mRNA respectively, through base pairing. U11 snRNP binds to a tandem repeat known as U11 snRNP-binding splicing enhancer (USSE) and initiates the splicing process. Since both U11 and U12 snRNAs come together as a bicomplex, they form a molecular bridge between two ends of introns in the pre-spliceosomal complex. The U11-48K and U11/U12-65K proteins recognize the splice site of U12 type intron and stabilize the U11/U12 bi-complex. After activating the spliceosomal complex, U11 snRNA leaves the assembly.
Role during alternative splicing:
This kind of 5' splice site recognition and intron bridging through protein-protein, protein-RNA and RNA-RNA interactions is unique in the minor splicesomal complex, unlike the major spliceosomal one. Since alternative splicing is the key to the variation of gene expression (mRNA) encoding proteins, U11 is crucial to this regulatory process and responsible in forming a proteomic pool. Therefore U11 snRNA is important in terms of evolutionary aspects. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Winemaking in Bordeaux**
Winemaking in Bordeaux:
Winemaking in Bordeaux is the collection of processes used to make Bordeaux wine. Winemaking in Bordeaux differs from winemaking in other places both because of the Bordelais climate and because of the particular Bordeaux style which the wine-maker is aiming for. A large number of processes are involved: pruning, training, spraying, (green harvesting), harvesting, (sorting), (de-stemming), crushing, fermentation, pressing, barrel-ageing, blending, bottling and bottle-ageing.
Viticulture:
Pests Bordeaux is a relatively humid region. Thus it is a place rife with diseases and other problems that afflict vines, compared with many of the world's other wine regions, such as Chile or Australia. oidium, mildew, coulure (failure of the flowers), millerandage (irregular ripening of the grapes), Eutypiose, Esca, Vers de la grappe and Botrytis (can be beneficial—see Sauternes) are the most common diseases or problems that occur.
Viticulture:
Pruning In Bordeaux, the pruning of the vine happens almost always as cane-pruning (as opposed to spur-pruning). There are two types of cane-pruning: guyot simple and guyot double. The simple way is frequently seen on the right bank, double most often on the left. Related to pruning is the trellising, where vines are dispersed along wires. It has become increasingly popular to raise the height of the trellis to the benefit of the grapes but to the discomfort for the vigneron.
Viticulture:
Use of fertilizer and other chemicals The use of chemical sprays and fertilizers has dropped in the recent decades in Bordeaux. Forty years ago, using fertilizers and different herbicides and fungicides were common, and made work easier for the manager. It also lowered the quality of the grapes, however. This use is still taking place in Bordeaux—but less and less so. Fertilizers, if used at all, are now more commonly supplied by compost, rather than man-made chemicals. Ploughing has replaced many pesticides and de-leafing has replaced fungicide use. While a healthier approach to agriculture has certainly come to Bordeaux, the châteaux have not adopted the biodynamic trend so popular in other wine regions, though Bordeaux is not entirely unfamiliar with the concept as 25 % of the vines arevdedicated to certified organic cultivation without chemical fertilisers and pesticides. Instead, the lutte raisonnée method is gaining ground.
Viticulture:
Green harvesting Bordeaux has seen a rise in the use of green harvesting, where unripe bunches are cut off in the summer in order to channel more of the plant's strength to the remaining bunches. While it is a popular process, it also has its opponents, such as Jean Gautreau of Château Sociando-Mallet, Gonzague Lurton of Château Durfort-Vivens and Olivier Bernard of Domaine de Chevalier, who claim that the remaining berries simply grow bigger, not better. Green harvesting requires cheap labor, often ignorant of properly cutting vines. Opponents point to vintages 1929 and 1947, which were high-yield and of great quality—and made entirely without green harvest techniques.
Viticulture:
Harvesting When harvest time approaches the Bordeaux wine producers start getting anxious. Unlike many other wine regions, weather in Bordeaux is relatively unstable and sudden changes in weather can delay a harvest, force a harvest in bad weather (diluting the wine) or severely damage the harvest. The appellations around Sauternes are even more vulnerable as certain micro-climatic conditions has to arrive plus they are forced to harvest late, risking the entire harvest to bad weather. Today the Bordelais Châteaux focus increasingly on the right time of harvest related to ripeness of the grapes. Cabernet grapes don't mature at the same rate as Merlot and thus picking both grapes at the same time rarely makes sense if optimal ripeness is sought. Thierry Manoncourt (Château Figeac) recollects: "In the past the whole vineyard would have been picked in eight days. Today it takes us twenty to thirty days."In Bordeaux, hand picking is now common among the more prestigious châteaux. But while hand-picking is foremost, some classified châteaux still harvest by machine. Mechanical harvesting also has its advantages, such as flexibility: it makes possible harvesting at night, which is preferable during hot weather. While the harvesting machines today have advanced in technology making them still more attractive, the delicate and selective process of harvesting by hand is still the best way to secure a maximum quality harvest. One problem with manual harvesting is the sheer size of vineyards in Bordeaux (not to mention the labor cost of hand-picking), with tens of thousands of hectares needing harvesting within a few weeks. The flatter geography of Bordeaux also allows for mechanical harvesting, whereas the steep slopes of wine-producing areas such as Côte-Rôtie make machine harvesting nearly impossible.
Viticulture:
Yields Yields in Bordeaux, as is the case in all other French AOCs, are capped by administrative rules. In Bordeaux this cap is 60 hl/ha with the option to raise yields by 20% by permission from the INAO. Any excess wine is sent for distillation. The yield is essential for the quality of the wine along with many other factors such as terroir. When making wine from a mediocre terroir, the producer has to lower his yields more than if it came from a superior terroir if he wants the wine to be of similar quality. In Bordeaux the recent decades have seen more and more focus on low yields. In the 1980s it was common, even for prestigious châteaux, to harvest the legal maximum. The reason is clear: if you have a good (expensive) label it is tempting to harvest 40, 50 or 60 hl/ha rather than 30 hl/ha as it means more money. But hl/ha isn't everything: as Michel Cazes of Château Lynch-Bages says: "When people talk about yields they forget about density. Here in the Médoc we have 10,000 vines per hectare. The crop expressed in hectolitres tells you nothing. Here in Pauillac and St-Julien and St-Estèphe I am sure that fifty to sixty is about right. Latour always had some of the highest yields in the region, but that was because none of their vines were missing"
Winemaking:
In Bordeaux, almost all wines are blended. Only a few producers make single-variety or varietal wines, though the lack of naming grape varieties on labels masks the fact. The typical blend consists of Merlot and Cabernet Sauvignon and/or Cabernet Franc, with small additions of Petit Verdot and Malbec and very occasionally Carmenere. Merlot is favored on the right bank of the Gironde River system, and Cabernet Sauvignon on the left, though Merlot acreage has been increasing on the left bank over the last decade or two. Today, winemaking in Bordeaux is a highly controlled process, with widespread use of stainless steel vats for fermentation, cooling apparatus, and a high degree of hygienic discipline. In 1951, chaptalization (adding sugar) became legal (it had likely taken place illegally prior to 1951). The use of chaptalization is common in Bordeaux, except in the warmest of vintages, and especially on the left bank, where Cabernet Sauvignon dominates and ripens later than Merlot.
Winemaking:
Grape processing Today, sorting and de-stemming are common techniques in Bordeaux and have been for some time. Great efforts have been made to improve these processes. Technology has also affected the crushing of the grapes, which had been done by treading since ancient times. More recently, machines have made crushing cheaper and safer, but they are less gentle with the grapes—breaking the pips releases unwanted tannins into the must. Today, some châteaux, such as Château Smith-Haut-Lafitte, do not crush the grapes at all, letting the fermentation begin within each grape (a process widely used in the Beaujolais region).
Winemaking:
To move the grapes, a number of wineries have stopped using pumps. Instead, after the crushing, they raise the grapes by conveyor belt. This is a gentler process, using gravity, rather than a pumping system.
Winemaking:
Fermentation Fermentation usually takes place in stainless steel vats, a technique introduced in the 1960s (lined cement vats were introduced already in the 1920s), to improve hygiene and control over the fermentation process (especially of temperature). During the 1980s, some producers began reintroducing wooden fermentation vats. There are pros and cons with all types of vats, and their role in winemaking seems less important than other elements in the process.Use of concentrators, where a winemaker can remove water from the must, is common in Bordeaux. Some producers (Christian Moueix of Pétrus, Anthony Barton of Château Leoville-Barton, Philippe Dhalluin of Château Mouton-Rothschild) are opposed to concentration, although others (such as Château Pomeaux) are big fans. While this process can certainly improve a wine in mediocre years, it is also open to abuse—with the result being an over-concentrated and poorly balanced wine.
Winemaking:
Pressing After fermentation comes the pressing. Bordeaux, along with other regions, has switched from hydraulic presses to the pneumatic press, where a pneumatic bladder filled with air results in a more gentle pressing of the wine. A third type of press is the vertical or hydraulic press. This is the most traditional, and also a gentle, type of press. However, is a very labour-intensive process.
Winemaking:
Oxygenation The modern, and very popular, method of micro-oxygenation, where microscopic amounts of oxygen are added to the wine during fermentation to stabilize (green) tannins and anthocyanins, has also caught on in Bordeaux. The most prestigious châteaux avoid the procedure, preferring to harvest grapes without green tannins. Micro-oxygenation is also used later in the process, during élevage, as a way of avoiding racking and controlling the amount of oxygen applied to the wine. (Racking allows for no such control). In this stage, however, the prestigious châteaux have fewer reservations, although not all producers are fans of micro-oxygenation during élevage.
Winemaking:
Aging In Bordeaux, most serious wines undergo barrel-ageing, although white wines can be an exception. Usually, six months of ageing in-barrel is required, but some châteaux barrel-age for as much as 20 months. The number of new barrels (which impart a higher degree of oak flavor to the wine) can vary from vintage to vintage, just as the duration of barrel-ageing. Only recently, addition of oak chips has been made legal in Bordeaux. During barrel-ageing, the wine needs to be racked in order to clear it of lees. This process is being challenged by some producers, as mentioned abovem since ageing on the lees can also add richness to the wine.
Winemaking:
Blending Once the producer decides the wine has aged for the right amount of time, the selection begins. The winemaker (or his/her team) find the right blend for the vintage. This is released as the château's grand vin. Inevitably, there will be some wine left—either of inferior quality or leftovers from the blending. This is usually released as a second-wine (or in some cases even a third-wine). While in theory inferior wine, some châteaux second-wine is of superior quality to other châteaux' grand vin and fetches high prices. Increasing the amount of second-wine can be a very conscious decision on the part of a winemaker, as a way of making a more and more superior grand vin - able to compete with the most prestigious wines in tastings.After blending, the wine will be bottled, and will then usually undergo a further ageing process before being released.
Winemaking:
In Bordeaux the oenologists play a huge role. Many oenologists work as consultants to different châteaux and carry much weight in major decisions regarding the wine. Amongst the most famous oenologists are Emile Peynaud, Jacques Boissenot, Pascal Chantonnet, Olivier Dauga, Stéphane Derenoncourt, Denis Dubourdieu, Jean-Philippe Fort, Gilles Pauquet, Michel Rolland, Stéphane Toutoundji and Christian Veyry. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Field Deployable Hydrolysis System**
Field Deployable Hydrolysis System:
The Field Deployable Hydrolysis System (FDHS) is a transportable, high throughput neutralization system developed by the U.S. Army for converting chemical warfare material into compounds not usable as weapons.
Operation:
Neutralization is facilitated through chemical reactions involving reagents that are mixed and heated to increase destruction efficiency, which is rated at 99.9 percent.The transportable FDHS is a self-contained system that includes power generators and a laboratory. Operational inputs include consumable materials such as water, reagents and fuel. It is designed to be set up within 10 days and is equipped with redundant critical systems. An on-site a crew of 15 trained personnel, including SME support, is needed for each shift of a possible 24-hour operational cycle.
Development:
A 20-week design and development phase was funded by the Defense Threat Reduction Agency in February 2013. The effort to develop a functional prototype was led by subject-matter experts from the Edgewood Chemical Biological Center (ECBC) in partnership with the United States Army Chemical Materials Agency. An operational model was developed over the course of six months, with the participation of 50 ECBC employees.
Deployment:
Two of these units were deployed on the MV Cape Ray (T-AKR-9679) for use in the destruction of Syria's chemical weapons. They are the "centerpiece" of the disarmament effort. The United Kingdom gave the United States £2.5 million of specialist equipment and training to enable the highest-priority chemicals to be processed more quickly. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**POMGNT1**
POMGNT1:
Protein O-linked-mannose beta-1,2-N-acetylglucosaminyltransferase 1 is an enzyme that in humans is encoded by the POMGNT1 gene.
Function and expression:
The product of the POMGNT1 gene, protein O-mannose beta-1,2-N-acetylglucosaminyltransferase (POMGnT1), participates in O-mannosyl glycan synthesis. A mutation in this gene is the cause of muscle-eye-brain disease (MIM 253280).Transcription of the POMGNT1 gene gives rise to a 2.7 kb mRNA in different tissues, with higher expression levels in the skeletal muscle, heart, and kidney and lower levels in the brain. POMGnT1 (EC 2.4.1.101) is a protein belonging to the GT13 family of glycosyltransferases according to the Carbohydrate-Active enZYmes (CAZy) database. In humans, the main isoform of POMGnT1 contains 660 amino acids whose sequence yields a calculated molecular mass of 75,252 Da (UniProtKB Q8WZA1).
Function and expression:
The POMGNT1 mRNA and its encoded protein is expressed in the neural retina of all mammals studied. POMGnT1 locates in the cytoplasmic fraction in the mouse retina, where it concentrates in the Golgi complex within the myoid of photoreceptor inner segments. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Silicide carbide**
Silicide carbide:
Silicide carbides or carbide silicides are compounds containing anions composed of silicide (Si4−) and carbide (C4−) or clusters therof. They can be considered as mixed anion compounds or intermetallic compounds, as silicon could be considered as a semimetal. Related compounds include the germanide carbides, phosphide silicides, boride carbides and nitride carbides. Other related compounds may contain more condensed anion combinations such as the carbidonitridosilicates with C(SiN3)4 with N bridging between two silicon atoms.
Production:
Silicide carbide compounds can be made by heating silicon, graphite, and metal together. It is important to exclude oxygen before and during the reaction. The flux method involves a reaction in a molten metal. Gallium is suitable, because it dissolves carbon and silicon, but does not react with them.
Properties:
Silicide carbides are a kind of ceramic, yet they also have metallic properties. They are not as brittle as most ceramics, but are stiffer than metals. They have high melting temperatures.In air silicide carbide compounds are stable, and are hardly affected by water. The appearance is often metallic grey. When powdered the colour is dark grey.When ErFe2SiC is dissolved in acid, mostly methane is produced, but the products include some hydrocarbons with two and three carbon atoms.The lanthanide contraction is evident with the cell sizes for rare earth element silicide carbides. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**International Journal of Computational Intelligence and Applications**
International Journal of Computational Intelligence and Applications:
The International Journal of Computational Intelligence and Applications is a refereed scientific journal published by Imperial College Press since 2001. It covers the theory and applications of computational intelligence and aims to provide "a vehicle whereby ideas using two or more conventional and computational intelligence based techniques could be discussed", such as neuro-fuzzy or evolutionary-symbolic combinations.
Abstracting and indexing:
The journal is abstracted and indexed in CompuScience and Inspec. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Distorted Schwarzschild metric**
Distorted Schwarzschild metric:
In physics, the distorted Schwarzschild metric is the metric of a standard/isolated Schwarzschild spacetime exposed in external fields. In numerical simulation, the Schwarzschild metric can be distorted by almost arbitrary kinds of external energy–momentum distribution. However, in exact analysis, the mature method to distort the standard Schwarzschild metric is restricted to the framework of Weyl metrics.
Standard Schwarzschild as a vacuum Weyl metric:
All static axisymmetric solutions of the Einstein–Maxwell equations can be written in the form of Weyl's metric, (1)ds2=−e2ψ(ρ,z)dt2+e2γ(ρ,z)−2ψ(ρ,z)(dρ2+dz2)+e−2ψ(ρ,z)ρ2dϕ2, From the Weyl perspective, the metric potentials generating the standard Schwarzschild solution are given by ln ln L2−M2l+l−, where (3)L=12(l++l−),l+=ρ2+(z+M)2,l−=ρ2+(z−M)2, which yields the Schwarzschild metric in Weyl's canonical coordinates that (4)ds2=−L−ML+Mdt2+(L+M)2l+l−(dρ2+dz2)+L+ML−Mρ2dϕ2.
Weyl-distortion of Schwarzschild's metric:
Vacuum Weyl spacetimes (such as Schwarzschild) respect the following field equations, 5.
a)∇2ψ=0, 5.
b)γ,ρ=ρ(ψ,ρ2−ψ,z2), 5.
c)γ,z=2ρψ,ρψ,z, 5.
d)γ,ρρ+γ,zz=−(ψ,ρ2+ψ,z2), where := ∂ρρ+1ρ∂ρ+∂zz is the Laplace operator.
Derivation of vacuum field equations. The vacuum Einstein's equation reads Rab=0 , which yields Eqs(5.a)-(5.c).
Moreover, the supplementary relation R=0 implies Eq(5.d). End derivation.
Eq(5.a) is the linear Laplace's equation; that is to say, linear combinations of given solutions are still its solutions. Given two solutions {ψ⟨1⟩,ψ⟨2⟩} to Eq(5.a), one can construct a new solution via (6)ψ~=ψ⟨1⟩+ψ⟨2⟩, and the other metric potential can be obtained by (7)γ~=γ⟨1⟩+γ⟨2⟩+2∫ρ{(ψ,ρ⟨1⟩ψ,ρ⟨2⟩−ψ,z⟨1⟩ψ,z⟨2⟩)dρ+(ψ,ρ⟨1⟩ψ,z⟨2⟩+ψ,z⟨1⟩ψ,ρ⟨2⟩)dz}.
Let ψ⟨1⟩=ψSS and γ⟨1⟩=γSS , while ψ⟨2⟩=ψ and γ⟨2⟩=γ refer to a second set of Weyl metric potentials. Then, {ψ~,γ~} constructed via Eqs(6)(7) leads to the superposed Schwarzschild-Weyl metric (8)ds2=−e2ψ(ρ,z)L−ML+Mdt2+e2γ(ρ,z)−2ψ(ρ,z)(L+M)2l+l−(dρ2+dz2)+e−2ψ(ρ,z)L+ML−Mρ2dϕ2.
With the transformations cos cos θ, sin cos 2θ, one can obtain the superposed Schwarzschild metric in the usual {t,r,θ,ϕ} coordinates, 10 sin 2θdϕ2.
The superposed metric Eq(10) can be regarded as the standard Schwarzschild metric distorted by external Weyl sources. In the absence of distortion potential {ψ(ρ,z)=0,γ(ρ,z)=0} , Eq(10) reduces to the standard Schwarzschild metric 11 sin 2θdϕ2.
Weyl-distorted Schwarzschild solution in spherical coordinates:
Similar to the exact vacuum solutions to Weyl's metric in spherical coordinates, we also have series solutions to Eq(10). The distortion potential ψ(r,θ) in Eq(10) is given by the multipole expansion 12 cos θ)M)Pi with := cos 2θ]1/2 where 13 := cos θR) denotes the Legendre polynomials and ai are multipole coefficients. The other potential γ(r,θ) is 14 )γ(r,θ)=∑i=1∞∑j=0∞aiaj (iji+j) (RM)i+j (PiPj−Pi−1Pj−1) −1M∑i=1∞αi∑j=0i−1 cos cos θ)] (RM)jPj. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Sibernetic**
Sibernetic:
Sibernetic is a fluid mechanics simulator developed for simulations of C. elegans in the OpenWorm project developed for the OpenWorm project by Andrey Palyanov, Sergey Khayrulin and Mike Vella as part of the OpenWorm team. Sibernetic provides an implementation of the PCISPH contractile matter algorithm for simulating muscle tissue and is applied to C. elegans locomotion. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Emergent virus**
Emergent virus:
An emergent virus (or emerging virus) is a virus that is either newly appeared, notably increasing in incidence/geographic range or has the potential to increase in the near future. Emergent viruses are a leading cause of emerging infectious diseases and raise public health challenges globally, given their potential to cause outbreaks of disease which can lead to epidemics and pandemics. As well as causing disease, emergent viruses can also have severe economic implications. Recent examples include the SARS-related coronaviruses, which have caused the 2002-2004 outbreak of SARS (SARS-CoV-1) and the 2019–21 pandemic of COVID-19 (SARS-CoV-2). Other examples include the human immunodeficiency virus which causes HIV/AIDS; the viruses responsible for Ebola; the H5N1 influenza virus responsible for avian flu; and H1N1/09, which caused the 2009 swine flu pandemic (an earlier emergent strain of H1N1 caused the 1918 Spanish flu pandemic). Viral emergence in humans is often a consequence of zoonosis, which involves a cross-species jump of a viral disease into humans from other animals. As zoonotic viruses exist in animal reservoirs, they are much more difficult to eradicate and can therefore establish persistent infections in human populations.Emergent viruses should not be confused with re-emerging viruses or newly detected viruses. A re-emerging virus is generally considered to be a previously appeared virus that is experiencing a resurgence, for example measles. A newly detected virus is a previously unrecognized virus that had been circulating in the species as endemic or epidemic infections. Newly detected viruses may have escaped classification because they left no distinctive clues, and/or could not be isolated or propagated in cell culture. Examples include human rhinovirus (a leading cause of common colds which was first identified in 1956), hepatitis C (eventually identified in 1989), and human metapneumovirus (first described in 2001, but thought to have been circulating since the 19th century). As the detection of such viruses is technology driven, the number reported is likely to expand.
Zoonosis:
Given the rarity of spontaneous development of new virus species, the most frequent cause of emergent viruses in humans is zoonosis. This phenomenon is estimated to account for 73% of all emerging or re-emerging pathogens, with viruses playing a disproportionately large role. RNA viruses are particularly frequent, accounting for 37% of emerging and re-emerging pathogens. A broad range of animals - including wild birds, rodents and bats - are associated with zoonotic viruses. It is not possible to predict specific zoonotic events that may be associated with a particular animal reservoir at any given time.Zoonotic spillover can either result in self-limited 'dead-end' infections, in which no further human-human transmission occurs (as with the rabies virus), or in infectious cases, in which the zoonotic pathogen is able to sustain human-human transmission (as with the Ebola virus). If the zoonotic virus is able to maintain successful human-human transmission, an outbreak may occur. Some spillover events can also result in the virus adapting exclusively for human infection (as occurred with the HIV virus), in which case humans become a new reservoir for the pathogen.
Zoonosis:
A successful zoonotic 'jump' depends on human contact with an animal harbouring a virus variant that is able to infect humans. In order to overcome host-range restrictions and sustain efficient human-human transmission, viruses originating from an animal reservoir will normally undergo mutation, genetic recombination and reassortment. Due to their rapid replication and high mutation rates, RNA viruses are more likely to successfully adapt for invasion of a new host population.
Zoonosis:
Examples of animal sources Bats While bats are essential members of many ecosystems, they are also frequently implicated as frequent sources of emerging virus infections. Their immune systems have evolved in such a way as to suppress any inflammatory response to viral infections, thereby allowing them to become tolerant hosts for evolving viruses, and consequently provide major reservoirs of zoonotic viruses. They are associated with more zoonotic viruses per host species than any other mammal, and molecular studies have demonstrated that they are the natural hosts for several high-profile zoonotic viruses, including severe acute respiratory syndrome-related coronaviruses and Ebola/Marburg hemorrhagic fever filoviruses. In terms of their potential for spillover events, bats have taken over the leading role previously assigned to rodents. Viruses can be transmitted from bats via several mechanisms, including bites, aerosolization of saliva (e.g. during echolocation), and faeces/urine.Due to their distinct ecology/behaviour, bats are naturally more susceptible to viral infection and transmission. Several bat species (e.g. brown bats) aggregate in crowded roosts, which promotes intra- and interspecies viral transmission. Moreover, as bats are widespread in urban areas, humans occasionally encroach on their habitats which are contaminated with guano and urine. Their ability to fly and migration patterns also means that bats are able to spread disease over a large geographic area, while also acquiring new viruses. Additionally, bats experience persistent viral infections which, together with their extreme longevity (some bat species have lifespans of 35 years), helps to maintain viruses and transmit them to other species. Other bat characteristics which contribute to their potency as viral hosts include: their food choices, torpor/hibernation habits, and susceptibility to reinfection.
Drivers of viral emergence:
Viral emergence is often a consequence of both nature and human activity. In particular, ecological changes can greatly facilitate the emergence and re-emergence of zoonotic viruses. Factors such as deforestation, reforestation, habitat fragmentation and irrigation can all impact the ways in which humans come into contact with wild animal species, and consequently promote virus emergence. In particular, habitat loss of reservoir host species plays a significant role in emerging zoonoses. Additionally, climate change can affect ecosystems and vector distribution, which in turn can affect the emergence of vector-borne viruses. Other ecological changes - for example, species introduction and predator loss - can also affect virus emergence and prevalence. Some agricultural practices, for example livestock intensification and inappropriate management/disposal of farm animal faeces, are also associated with an increased risk of zoonosis.Viruses may also emerge due to the establishment of human populations that are vulnerable to infection. For example, a virus may emerge following loss of cross-protective immunity, which may occur due to loss of a wild virus or termination of vaccination programmes. Well-developed countries also have higher proportions of aging citizens and obesity-related disease, thus meaning that their populations may be more immunosuppressed and therefore at risk of infection. Contrastingly, poorer nations may have immunocompromised populations due to malnutrition or chronic infection; these countries are also unlikely to have stable vaccination programmes. Additionally, changes in human demographics – for example, the birth and/or migration of immunologically naïve individuals – can lead to the development of a susceptible population that enables large-scale virus infection.
Drivers of viral emergence:
Other factors which can promote viral emergence include globalisation; in particular, international trade and human travel/migration can result in the introduction of viruses into new areas. Moreover, as densely populated cities promote rapid pathogen transmission, uncontrolled urbanization (i.e. the increased movement and settling of individuals in urban areas) can promote viral emergence. Animal migration can also lead to the emergence of viruses, as was the case for the West Nile virus which was spread by migrating bird populations. Additionally, human practices regarding food production and consumption can also contribute to the risk of viral emergence. In particular, wet markets (i.e. live animal markets) are an ideal environment for virus transfer, due to the high density of people and wild/farmed animals present. Consumption of bushmeat is also associated with pathogen emergence.
Prevention:
Control and prevention of zoonotic diseases depends on appropriate global surveillance at various levels, including identification of novel pathogens, public health surveillance (including serological surveys), and analysis of the risks of transmission. The complexity of zoonotic events around the world predicates a multidisciplinary approach to prevention. The One Health Model has been proposed as a global strategy to help prevent the emergence of zoonotic diseases in humans, including novel viral diseases. The One Health concept aims to promote the health of animals, humans, and the environment, both locally and globally, by fostering understanding and collaboration between practitioners of different interrelated disciplines, including wildlife biology, veterinary science, medicine, agriculture, ecology, microbiology, epidemiology, and biomedical engineering.
Virulence of emergent viruses:
As hosts are immunologically naïve to pathogens they have not encountered before, emergent viruses are often extremely virulent in terms of their capacity to cause disease. Their high virulence is also due to a lack of adaptation to the new host; viruses normally exert strong selection pressure on the immune systems of their natural hosts, which in turn exerts a strong selection pressure on viruses. This coevolution means that the natural host is able to manage infection. However, when the virus jumps to a new host (e.g. humans), the new host is unable to deal with infection due to a lack of coevolution, which results in mismatch between host immunoeffectors and virus immunomodulators.Additionally, in order to maximise transmission, viruses often naturally undergo attenuation (i.e. virulence is reduced) so that infected animals can survive long enough to infect other animals more efficiently. However, as attenuation takes time to achieve, new host populations will not initially benefit from this phenomenon. Moreover, as zoonotic viruses also naturally exist in animal reservoirs, their survival is not dependent on transmission between new hosts; this means that emergent viruses are even more unlikely to attenuate for the purpose of maximal transmission, and they remain virulent.Although emergent viruses are frequently highly virulent, they are limited by several host factors including: innate immunity, natural antibodies and receptor specificity. If the host has previously been infected by a pathogen that is similar to the emergent virus, the host may also benefit from cross-protective immunity.
Examples of emergent viruses:
Influenza A Influenza is a highly contagious respiratory infection, which affects approximately 9% of the global population and causes 300,000 to 500,000 deaths annually. Based on their core proteins, influenza viruses are classified into types A, B, C and D. While both influenza A and B can cause epidemics in humans, influenza A also has pandemic potential and a higher mutation rate, therefore is most significant to public health.Influenza A viruses are further classified into subtypes, based on the combinations of the surface glycoproteins hemagglutinin (HA) and neuraminidase (NA). The primary natural reservoir for most influenza A subtypes are wild aquatic birds; however, through a series of mutations, a small subset of these viruses have adapted for infection of humans (and other animals). A key determinant of whether a particular influenza A subtype can infect humans is its binding specificity. Avian influenza A preferentially binds to cell surface receptors with a terminal α2,3‐linked sialic acid, while human influenza A preferentially binds to cell surface receptors with a terminal α2,6‐linked sialic acid. Via mutation, some avian influenza A viruses have successfully altered their binding specificity from α2,3‐ to α2,6‐linked sialic acid. However, in order to emerge in humans, avian influenza A viruses must also adapt their RNA polymerases for function in mammalian cells, as well as mutating for stability in the acidic respiratory tract of humans.Following adaptation and host switch, influenza A viruses have the potential to cause epidemics and pandemics in humans. Minor changes in HA and NA structure (antigenic drift) occur frequently, which enables the virus to cause repetitive outbreaks (i.e. seasonal influenza) by evading immune recognition. Major changes in HA and NA structure (antigenic shift), which are caused by genetic reassortment between different influenza A subtypes (e.g. between human and animal subtypes), can instead cause large regional/global pandemics. Due to the emergence of antigenically different influenza A strains in humans, four pandemics occurred in the 20th century alone.Additionally, although animal influenza A viruses (e.g. swine influenza) are distinct from human influenza viruses, they can still cause zoonotic infection in humans. These infections are largely acquired following direct contact with infected animals or contaminated environments, but do not result in efficient human-human transmission; examples of this include H5N1 influenza and H7N9 influenza.
Examples of emergent viruses:
SARS-CoV In 2002, a highly pathogenic SARS-CoV (Severe Acute Respiratory Syndrome Coronavirus) strain emerged from a zoonotic reservoir; approximately 8000 people were infected worldwide, and mortality rates approached 50% or more in the elderly. As SARS-CoV is most contagious post-symptoms, the introduction of strict public health measures effectively halted the pandemic. The natural reservoir host for SARS-CoV is thought to be horseshoe bats, although the virus has also been identified in several small carnivores (e.g. palm civets and racoon dogs). The emergence of SARS-CoV is believed to have been facilitated by Chinese wet markets, in which civets positive for the virus acted as intermediate hosts and passed SARS-CoV onto humans (and other species). However, more recent analysis suggests that SARS-CoV may have directly jumped from bats to humans, with subsequent cross-transmission between humans and civets.In order to infect cells, SARS-CoV uses the spike surface glycoprotein to recognise and bind to host ACE-2, which it uses as a cellular entry receptor; the development of this characteristic was crucial in enabling SARS-CoV to ‘jump’ from bats to other species.
Examples of emergent viruses:
MERS-CoV First reported in 2012, MERS-CoV (Middle East Respiratory Syndrome Coronavirus) marks the second known introduction of a highly pathogenic coronavirus from a zoonotic reservoir into humans. The case mortality rate of this emergent virus is approximately 35%, with 80% of all cases reported by Saudi Arabia. Although MERS-CoV is likely to have originated in bats, dromedary camels have been implicated as probable intermediate hosts. MERS-CoV is believed to have been circulating in these mammals for over 20 years, and it is thought that novel camel farming practices drove the spillover of MERS-CoV into humans. Studies have shown that humans can be infected with MERS-CoV via direct or indirect contact within infected dromedary camels, while human-human transmission is limited.MERS-CoV gains cellular entry by using a spike surface protein to bind to the host DPP4 surface receptor; the core subdomain of this spike surface protein shares similarities with that of SARS-CoV, but its receptor binding subdomain (RBSD) significantly differs.
Examples of emergent viruses:
Bluetongue disease Bluetongue disease is a non-contagious vector-borne disease caused by bluetongue virus, which affects species of ruminants (particularly sheep). Climate change has been implicated in the emergence and global spread of this disease, due to its impact on vector distribution. The natural vector of the bluetongue virus is the African midge C. imicola, which is normally limited to Africa and subtropical Asia. However, global warming has extended the geographic range of C. imicola, so that it now overlaps with a different vector (C. pulcaris or C. obsoletus) with a much more northward geographic range. This change enabled the bluetongue virus to jump vector, thus causing the northward spread of bluetongue disease into Europe. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Kinin**
Kinin:
A kinin is any of various structurally related polypeptides, such as bradykinin and kallidin. They are members of the autacoid family. Kinins are peptides that are cleaved from kininogens by the process of kallikreins. Kallikreins activate kinins when stimulated.It is a component of the kinin-kallikrein system.
Their precursors are kininogens. Kininogens contain a 9-11 amino acid bradykinin sequence.In botany, the plant hormones known as cytokinins were first called kinins, but the name was changed to avoid confusion.
Effects of Kinins:
Kinin are short lived peptides that cause pain sensation, arteriolar dilation, increase vascular permeability and cause contractions in smooth muscle. Kinins transmit their effects through G protein- coupled receptors.Kinin act on axons to block nervous impulses, which leads to distal muscle relaxation. Kinin are also potent nerve stimulators. which is mostly responsible for the sense of pain (and sometimes itching). Kinin increase vascular permeability by acting on vascular endothelial cells to cause cell contraction. Concomitantly they induce local expression of adhesive molecules. Together they increase leukocytes adhesion and extravasation. Kinin are rapidly inactivated by the proteases locally generated during the above mentioned processes.They act locally to induce vasodilation and contraction of smooth muscle. Kinins function as mediators for inflammatory responses by triggering the immune system. They are also able to regulate cardiovascular and renal function through mediating the effects of ACE inhibitors. Reduced kinin activity can result in high blood pressure, sodium retention and the narrowing of blood vessels.Aspirin inhibits the activation of kallenogen by interfering with the formation of kallikrein enzyme which is essential in the process of activation.
Where kinins are produced:
Kinins are mostly produced at inflamed or injured tissue of the body and human body fluids. Kinin peptides (kallidin and bradykinin) are located in human blood and urine.
Kinin Receptors:
There are two types of kinin receptors, B1 and B2. B1 receptors are G-protein coupled receptors that are induced by injured tissue. The quantity of B2 receptors in the human body, exceed B1 receptors.
History:
Kinin was initially discovered by J.E. Abelous and E. Bardier in 1909 when performing experiments utilizing human body fluids. Human body fluids such as urine was injected into dogs and it was observed that the urine caused a reduction in blood pressure. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Final state conjecture**
Final state conjecture:
The final state conjecture is that the end state of the universe will consist of black holes and gravitational radiation. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Saturation mutagenesis**
Saturation mutagenesis:
Site saturation mutagenesis (SSM), or simply site saturation, is a random mutagenesis technique used in protein engineering, in which a single codon or set of codons is substituted with all possible amino acids at the position. There are many variants of the site saturation technique, from paired site saturation (saturating two positions in every mutant in the library) to scanning site saturation (performing a site saturation at every site in the protein, resulting in a library of size [20 x (number of residues in the protein)] that contains every possible point mutant of the protein).
Method:
Saturation mutagenesis is commonly achieved by site-directed mutagenesis PCR with a randomised codon in the primers (e.g. SeSaM) or by artificial gene synthesis, with a mixture of synthesis nucleotides used at the codons to be randomised.Different degenerate codons can be used to encode sets of amino acids. Because some amino acids are encoded by more codons than others, the exact ratio of amino acids cannot be equal. Additionally, it is usual to use degenerate codons that minimise stop codons (which are generally not desired). Consequently, the fully randomised 'NNN' is not ideal, and alternative, more restricted degenerate codons are used. 'NNK' and 'NNS' have the benefit of encoding all 20 amino acids, but still encode a stop codon 3% of the time. Alternative codons such as 'NDT', 'DBK' avoid stop codons entirely, and encode a minimal set of amino acids that still encompass all the main biophysical types (anionic, cationic, aliphatic hydrophobic, aromatic hydrophobic, hydrophilic, small). In the case there is no restriction to use a single degenerate codon only, it is possible to reduce the bias considerably. Several computational tools were developed to allow high level of control over the degenerate codons and their corresponding amino acids.
Applications:
Saturation mutagenesis is commonly used to generate variants for directed evolution. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Biopharmaceutics Classification System**
Biopharmaceutics Classification System:
The Biopharmaceutics Classification System is a system to differentiate drugs on the basis of their solubility and permeability.This system restricts the prediction using the parameters solubility and intestinal permeability. The solubility classification is based on a United States Pharmacopoeia (USP) aperture. The intestinal permeability classification is based on a comparison to the intravenous injection. All those factors are highly important because 85% of the most sold drugs in the United States and Europe are orally administered.
BCS classes:
According to the Biopharmaceutical Classification System (BCS) drug substances are classified to four classes upon their solubility and permeability: Class I - high permeability, high solubility Example: metoprolol, paracetamol Those compounds are well absorbed and their absorption rate is usually higher than excretion.
Class II - high permeability, low solubility Example: glibenclamide, bicalutamide, ezetimibe, aceclofenac The bioavailability of those products is limited by their solvation rate. A correlation between the in vivo bioavailability and the in vitro solvation can be found.
Class III - low permeability, high solubility Example: cimetidine The absorption is limited by the permeation rate but the drug is solvated very fast. If the formulation does not change the permeability or gastro-intestinal duration time, then class I criteria can be applied.
Class IV - low permeability, low solubility Example: Bifonazole Those compounds have a poor bioavailability. Usually they are not well absorbed over the intestinal mucosa and a high variability is expected.
Definitions:
The drugs are classified in BCS on the basis of solubility, permeability, and dissolution.
Definitions:
Solubility class boundaries are based on the highest dose strength of an immediate release product. A drug is considered highly soluble when the highest dose strength is soluble in 250 ml or less of aqueous media over the pH range of 1 to 7.5. The volume estimate of 250 ml is derived from typical bioequivalence study protocols that prescribe administration of a drug product to fasting human volunteers with a glass of water.
Definitions:
Permeability class boundaries are based indirectly on the extent of absorption of a drug substance in humans and directly on the measurement of rates of mass transfer across human intestinal membrane. Alternatively non-human systems capable of predicting drug absorption in humans can be used (such as in-vitro culture methods). A drug substance is considered highly permeable when the extent of absorption in humans is determined to be 90% or more of the administered dose based on a mass-balance determination or in comparison to an intravenous dose.
Definitions:
For dissolution class boundaries, an immediate release product is considered rapidly dissolving when no less than 85% of the labeled amount of the drug substance dissolves within 15 minutes using USP Dissolution Apparatus 1 at 100 RPM or Apparatus 2 at 50 RPM in a volume of 900 ml or less in the following media: 0.1 M HCl or simulated gastric fluid or pH 4.5 buffer and pH 6.8 buffer or simulated intestinal fluid. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Science and technology in Pakistan**
Science and technology in Pakistan:
Science and technology is a growing field in Pakistan and has played an important role in the country's development since its founding. Pakistan has a large pool of scientists, engineers, doctors, and technicians assuming an active role in science and technology. The real growth in science in Pakistan occurred after the establishment of the Higher education Commission in 2002 which supported science in a big way and also became the major sponsor of the Pakistan Academy of Sciences under the leadership of Prof. Atta-ur-Rahman. The emphasis was placed on quality rather than numbers during this period. The quality measures introduced by Prof. Atta-ur-Rahman as Founding Chairman HEC included:1) All Ph.D. thesis were evaluated by eminent foreign scientists,2) All PhD theses and research papers were checked for plagiarism 3) Some 11,000 students were sent abroad to leading universities for PhD level training and absorbed on their return, 4) Appointments at faculty positions were linked to international stature of the applicants as judged from their international publications, patents and citations, and (5) Quality Enhancement Cells were established in all universities for the first time in the history of the country. (6) The minimum criteria for establishment of a new university were approved by the Cabinet and universities that did not meet this criteria were closed down. (7) The Model University Ordinance was approved (Appendix 3 in the reference) setting the governance parameters for new universities. (8) A list of fake higher education institutions was prepared and made public. (9) Quality Assurance Agency (QAA) was set up within the Higher Education Commission that established Quality Enhancement Cells (QECs) as its operational units in public and private-sector universities across the country. (10) The funding of universities was linked to excellence in teaching and research under a formula based funding mechanism that considered enrolment, subjects and quality of teaching and research. The first IT policy and implementation strategy was approved under the leadership of Prof. Atta-ur-Rahman, then Federal Minister of Science & technology, in August 2000 which laid the foundations of the development of this sector On the request of Prof. Atta-ur-Rahman, Intel initiated a nationwide programme to train school teachers in Information and Communication technologies in March 2002 which has led to the training of 220,000 school teachers in 70 districts and cities across Pakistan. A 15-year tax holiday was approved on the recommendation of Prof. Atta-ur-Rahman which has resulted in growth of IT business from $30 million in 2001 to over $3 billion. The Pakistan Austria University of Applied Engineering (Fachhochschule) has been established in Haripur Hazara under a Steering Committee Chaired by Prof. Atta-ur-Rahman in which students will get degrees from several Austrian universities. Pakistan's growth in scientific output can be seen from the fact that in 1990 Pakistan published 926 scholarly documents while in 2018 the number rose to 20548, a twenty times increase.In contrast India published 21443 scholarly documents in 1990 and the number rose to 171356 in 2018, an eight times increase. In 2018, 336 people per million were researchers in the R&D (Research and Development sector) in Pakistan compared to 256 people per million being researchers in India. The reforms begun by Prof. Atta-ur-Rahman FRS in 2003-2008 have continued over the subsequent decade and according to the Web of Science report, there was a 300% growth in research publications in 2019 over the decade, with 2019 marking the first year in which Pakistan was ranked above the world average in research. In 2019, Pakistan produced 300% more publications indexed in the Web of Science Core Collection than in 2010. In the decade of 2010-2019, more than half of Pakistan’s research was published in journals with Impact Factor. The global influence of Pakistan’s research is increasing as scientists in the country are publishing more in top quartile journals. The Category Normalized Citation Impact of Pakistan’s publications (which measures publications’ impact against their peers worldwide) has risen from 0.67 to 1.03. output. As of 2020, Pakistan has 85% teledensity with 183 million celllular, 98 million 3G/4G and 101 million broadband subscribers, due to the foundations laid by Prof. Atta-ur-Rahman of the IT and telecom industry during 2000-2008. In an analysis of scientific research productivity of Pakistan, in comparison to Brazil, Russia, India and China, Thomson Reuters has applauded the developments that have taken place as a result of the reforms introduced by Prof. Atta-ur-Rahman FRS, since Pakistan has emerged as the country with the highest increase in the percentage of highly cited papers in comparison to the "BRIC" countriesChemistry remains the strongest subject in the country with the International Center for Chemical and Biological Sciences playing the lead role with the largest postgraduate research program in the country having about 600 students enrolled for PhD. Physics (theoretical, nuclear, particle, laser, and quantum physics), material science, metallurgy (engineering), biology, and mathematics, are some of the other fields in which Pakistani scientists have contributed. From the 1960s and onwards, the Pakistani government made the development and advancement of science a national priority and showered top scientists with honours. While the government has made efforts to make science a part of national development, there have been criticisms of federal policies, such as the government's dissolution of the Higher Education Commission of Pakistan (HEC)— an administrative body that supervised research in science – in 2011. This attempted dissolution failed to materialise because of a Supreme Court of Pakistan decision on a petition filed by Prof. Atta-ur-Rahman, former Federal Minister of Science & Technology and former founding Chairman of the Higher Education Commission. Pakistani scientists have also won acclaim in mathematics and in several branches of physical science, notably theoretical and nuclear physics, chemistry, and astronomy. Professor Abdus Salam, a theoretical physicist won the Nobel Prize in Physics in 1979, being the first and only Pakistani to date to have received the honor. Prof. Atta-ur-Rahman an organic chemist was elected as Fellow of Royal Society (London) in 2006 in recognition of his contributions in the field of natural products thereby becoming the first scientist from the Islamic world to receive this honour for work carried out within an Islamic country. The contributions of Prof. Atta-ur-Rahman to uplift science and higher education in Pakistan were internationally acknowledged and a tribute paid to him in the world's leading science journal Nature that termed him as "a force of nature". In an analysis of scientific research productivity of Pakistan, in comparison to Brazil, Russia, India, and China, Thomson Reuters has applauded the developments that have taken place as a result of the reforms introduced by Prof. Atta-ur-Rahman FRS, since Pakistan has emerged as the country with the highest increase in the percentage of highly cited papers in comparison to the "BRIC" countries. In recognition of building strong bridges between science in Pakistan and China, Prof. Atta-ur-Rahman FRS received the highest national award of China, the "International Science and Technology Cooperation Award". His book on NMR spectroscopy published by Springer Verlag was translated into Japanese language and used for teaching courses on NMR spectroscopy in Japan. His book entitled "Stereoselective Synthesis in Organic Chemistry" published by Springer Verlag was described as a "monumental tome" by the Nobel Laureate Sir Derek Barton who wrote the Foreword to this book.Technology is highly developed in nuclear physics and explosives engineering, where the arms race with India convinced policymakers to set aside sufficient resources for research. Due to a programme directed by Munir Ahmad Khan and the Pakistan Atomic Energy Commission (PAEC), Pakistan is the seventh nation to have developed an atomic bomb, which the global intelligence community believes it had done by 1983 (see Kirana-I), nine years after India (see Pokhran-I). Pakistan first publicly tested its devices (see Chagai-I and Chagai-II) on 28 and 30 May 1998, two weeks after India carried out its own tests (See Pokhran-II).Space exploration was hastily developed, in 1990 Pakistan launched Badr-1 followed by Badr-II in 2001. Since the 1980s, the space programme dedicated itself to military technologies (Space weapons programme and Integrated missile systems), and maintains a strong programme developed for military applications.
Science and technology in Pakistan:
Pakistan is an associate member of CERN, one of the few countries to obtain that status. Pakistan was ranked 99th in the Global Innovation Index in 2021, up from 105th in 2019.During 2018-2019, the Government of Pakistan has formed a number of Task Forces to strengthen science and technology, information technology and knowledge economy. The task force formed in 2018 on "Technology Driven Knowledge Economy" is chaired by the Prime Minister Mr. Imran Khan and has Atta-ur-Rahman as its Vice Chairman The group has several important Federal Ministers as members including Ministers of Finance, Planning, Education, IT/Telecom, Science & Technology and chairman Higher Education Commission. The task force aims to promote research in important and emerging technology fields. Another important task force of the Prime Minister is that on science & technology with Atta-ur-Rahman as its chairman. As a result of the efforts of these Task Forces under the leadership of Prof. Atta-ur-Rahman FRS, a huge change has occurred in the Ministry of Science and Technology and the development budget of the Federal Ministry of Science and technology has been enhanced by over 600% due to the projects initiated by these Task Forces, allowing a large number of new important initiatives in the fields of materials engineering, genomics, industrial biotechnology, alternative energy, minerals, regenerative medicine, neuroscience, and artificial intelligence to be undertaken. Pakistan's first foreign engineering university (Pak Austria Fachhochschule) is a unique hybrid model involving a Fachhochschule half and a postgraduate research half, with a central technology park. With 8 foreign universities collaborating (3 Austrian and 5 Chinese), it has also started functioning under the supervision of a steering committee headed by Atta-ur-Rahman in Haripur, Hazara. A number of such foreign engineering universities are in the process of being established under the supervision of Prof. Atta-ur-Rahman FRS. These include one in Sialkot the foundation stone of which has already been laid by the Prime Minister of Pakistan, and another in the lands behind Prime Minister House, Islamabad
History:
The Scientific and Technological Research Division (S&TR) was established in 1964 for (i) coordination and implementation of national science and technology policy; (ii) promotion and coordination of research and utilization of the results of research; (iii) development, production and utilization of nuclear energy; and (iv) coordination of utilization of scientific and technological manpower. The Division was administratively responsible for the National Science Council, the Council of Scientific and Industrial Research, the Atomic Energy Commission and Space and Upper Atmospheric Research Committee.
History:
The Ministry of Science and Technology (MoS&T) has been functioning since 1972. It is the national focal point and enabling arm of Government of Pakistan for planning, coordinating and directing efforts; to initiate and launch scientific and technological programs and projects as per national agenda for sound and sustainable Science & Technology Research base for the socio-economic development.
History:
From the areas of industrial development to renewable energy and rural development, the Ministry suggests technological development for higher growth-rates and to improve standards of living. Its principal focus is on building Pakistan's technological competence and developing a larger pool of human resources to reverse brain drain, and for integrating the existing technological infrastructure for the strengthening of technology institutions, effective governance of S&TR and enhancing the capacity of indigenous innovation systems.
History:
Golden age of science The 1960s and the 1970s period is regarded as the initial rise of Pakistan's science, which gained an international reputation in the different science communities of the world. During this period, scientists contributed to the fields of, particularly, Natural Product Chemistry, theoretical, particle, mathematical, and nuclear physics, and other major and subfields of Chemistry and Physics. The research was preceded by such scientists as Riazuddin, Ishfaq Ahmad, Salimuzzaman Siddiqui, Atta-ur-Rahman and Samar Mubarakmand. However, the major growth in scientific output occurred after the establishment of the Higher Education Commission which was accompanied by a 60-fold increase in funding for science The real growth of science in Pakistan occurred under the leadership of Prof. Atta-ur-Rahman during 2000–2008 when he was the Federal Minister of Science & Technology and later Chairman of the Higher Education Commission (HEC) with the status of Federal Minister. The chairperson of the Senate Standing Committee on Education announced the first 6 years of HEC under Prof. Atta-ur-Rahman as "Pakistan's golden period". Thomson Reuters, in an independent assessment of Pakistan's progress in international publications, has acknowledged that in the last decade there has been a fourfold increase in international publications and a tenfold growth in highly cited papers, statistics that were better than the BRIC countries.The remarkable transformation of science and higher education under the leadership of Prof. Atta-ur-Rahman as Federal Minister of Science & Technology and later as Chairman of Higher Education Commission with status of a Federal Minister during the period 2000–2008 was applauded by many independent experts and he was called a "force of nature" in a review published in NatureDr. Abdus Salam, the first Pakistani winner of the Nobel Prize in Physics, was the father of physics research in Pakistan. Under the watchful direction of Salam, mathematicians and physicists tackled the greatest and outstanding problems in physics and mathematics. From 1960 to 1974, Salam was responsible for leading the research at its maximum point. This prompted the international recognition of Pakistani mathematicians and physicists, allowing them to conduct their research at CERN. Salam and his students (Riazuddin, Fayyazuddin, and others) revolutionized particle and theoretical physics, are thought to be modern pioneers of particle physics at all aspect of it. Pure research was undertaken in Quantum electrodynamics, Quantum field theory, protonic decay and major fields in physics, were pioneered by Pakistan's scientists. With the establishment of nuclear and neutron institutes in the country, Pakistan's mathematicians introduced complex mathematical applications to study and examine the behaviours of elements during the fission process. Salimuzzaman Siddiqui, Atta-ur-Rahman and Iqbal Choudhary are the pioneering personalities for studying the isolation of unique chemical compounds from the Neem (Azadirachta indica), Rauvolfia, periwinkle (Catharanthus roseus), (Buxus papillosa) and various other plants.
State controlled science:
Unlike some Western countries, the majority of the research programmes are conducted not at the institutions (such as universities) but at specially set up research facilities and institutes. These institutes are performed under the government's Ministry of Science that overlooks the development and promotion of science in the country, while others are performed under the Pakistan Academy of Sciences, other specialized academies and even the research arms of various government ministries. At first, the core of fundamental science was the Pakistan Academy of Sciences, originally set up in 1953 and moved from Karachi to Islamabad in 1964. The Pakistan Academy of Sciences has a large percentage of researchers in the natural sciences, particularly physics. From 1947 to 1971, the research was being conducted independently with no government influence. The High Tension Laboratories (HTL) at the Government College University, Lahore (GCU) was established by R. M. Chaudhry with funds given by the British government in the 1950s. In 1967, Professor Abdus Salam led the foundation of the Institute of Theoretical Physics (ITP) at the Quaid-e-Azam University, and the establishment of the Pakistan Institute of Nuclear Science and Technology (PINSTECH) and the Centre for Nuclear Studies; all were independently established by Pakistan's academic scientists with financial assistance provided by European countries. However, after Zulfikar Ali Bhutto became president, he took control of scientific research in 1972 as part of his intensified socialist reforms and policies. With advice taken from Dr. Mubashir Hassan, Bhutto established the Ministry of Science with Ishrat Hussain Usmani, a bureaucrat with a doctorate in atomic physics.During the 1950s and 1960s, both West Pakistan and East Pakistan had their own academies of science, with East Pakistan relying on West Pakistan to allot the funds. Medical research is coordinated and funded by the Health Ministry and agricultural research is led by Agriculture Ministry and likewise, the research on environmental sciences is headed by the Environment Ministry.The aftermath of the 1971 Indo-Pakistan Winter War was that President Bhutto increased scientific funding by the Government by more than 200%, mostly dedicated to military research and development. Bhutto, with the help of his Science Adviser Dr. Salam, gathered hundreds of Pakistani scientists working abroad to develop what became Pakistan's atom bomb. This crash programme was directed at first by Dr. Abdus Salam until 1974, and then directed and led by Dr. Munir Ahmad Khan from 1974 until 1991. For the first time, an effort was made by the government when Pakistan's citizens made advancements in nuclear physics, theoretical physics, and mathematics. In the 1980s, General Muhammad Zia-ul-Haq radicalized science by enforcing pseudoscience – by his Muslim fundamentalists as administrators – in Pakistan's schools and universities. Zia-ul-Haq later promoted Dr. Abdul Qadeer Khan to export the sensitive industrial (military) technologies to Libya, Iran, and North Korea. Because of government control, academic research in Pakistan remains highly classified and unknown to the international scientific community. There have been several failed attempts made by foreign powers to infiltrate the country's research facilities to learn how much research has progressed and how much clandestine knowledge has been gained by Pakistan's scientists. One of the notable cases was in the 1970s when the Libyan intelligence made an unsuccessful attempt to gain knowledge on critical aspects of nuclear technology, and crucial mathematical fast neutron calculations in theoretical physics. It was thwarted by the ISI Directorate for Joint Intelligence Technical (JIT). From the 1980s and onward, both Russian intelligence and the Central Intelligence Agency made several attempts to access Pakistan's research but because of the ISI, they were unable to gain any information. From the period 1980 to 2004, research in science fell short until General Pervez Mushrraf established the Higher Education Commission (HEC) which heightened the contribution of science and technology in Pakistan. The major boost to science in Pakistan occurred under the leadership of Prof. Atta-ur-Rahman as the founding Chairman of the Higher Education Commission when about 11,000 students were sent to top universities abroad for Ph.D. and postdoctoral training. This has resulted in the enormous increase in the research output of Pakistan in Impact factor journals from about 800 per year in the year 2000 to over 12,000 publications per year. This drew positive comments from Thomson Reuters about the sharp increase in highly cited papers in comparison to Brazil, Russia, India and China Major research was undertaken by Pakistan's institutes in the field of natural sciences. In 2003, the Ministry of Science and Technology of the Government of Pakistan and the United States Department of State signed a comprehensive Science and Technology Cooperation Agreement that established a framework to increase cooperation in science, technology, engineering and education for mutual benefit and peaceful purposes between the science and education communities in both countries. In 2005, the United States Agency for International Development (USAID) joined with the Ministry of Science and Technology (MOST) and the Higher Education Commission of Pakistan to support the joint Pakistan-U.S. Science and Technology Cooperation Program. Beginning in 2008, the U.S. Department of State joined USAID as U.S. co-sponsor of the program. This program, which is being implemented by the National Academy of Sciences on the U.S. side, is intended to increase the strength and breadth of cooperation and linkages between Pakistan scientists and institutions with counterparts in the United States. However, with unfavourable situations, research declined. In 2011, the government dissolved the HEC and the control of education was taken over by governmental ministries. Prof. Atta-ur-Rahman filed a petition in the Supreme Court of Pakistan against the government action. The Supreme Court decided in favour of the stand taken by Prof. Atta-ur-Rahman, and the federal nature of the Higher Education Commission was preserved.
Science policy:
National Science, Technology and Innovation Policy The Federal Ministry of Science and Technology has overseen the S&T sector since 1972. However, it was not until 2012 that Pakistan's first National Science, Technology and Innovation Policy was formulated: this was also the first time that the government had formally recognized innovation as being a long-term strategy for driving economic growth. The policy principally emphasizes the need for human resource development, endogenous technology development, technology transfer and greater international co-operation in research and development (R&D).The policy was informed by the technology foresight exercise undertaken by the Pakistan Council for Science and Technology from 2009 onwards. By 2014, studies had been completed in 11 areas: agriculture, energy, ICTs, education, industry, environment, health, biotechnology, water, nanotechnology, and electronics. Further foresight studies were planned on pharmaceuticals, microbiology, space technology, public health, sewage, and sanitation, as well as higher education.
Science policy:
National Science, Technology and Innovation Strategy Following the change of government in Islamabad after the May 2013 general election, the new Ministry of Science and Technology issued the draft National Science, Technology and Innovation Strategy 2014–2018, along with a request for comments from the public. This strategy has been mainstreamed into the government's long-term development plan, Vision 2025, a first for Pakistan.The central pillar of the draft National Science, Technology and Innovation Strategy is human development. Although the pathway to implementation is not detailed, the new strategy fixes a target of raising Pakistan's gross domestic expenditure on R&D (GERD) from 0.29% (2013) to 0.5% of GDP by 2015 then to 1% of GDP by the end of the current government's five-year term in 2018. The ambitious target of tripling the GERD/GDP ratio in just seven years is a commendable expression of the government's resolve but ambitious reforms will need to be implemented concurrently to achieve the desired outcome.
National prizes:
The most prestigious government prize awarded for achievements in science and technology is Nishan-e-Imtiaz (or in English Order of Excellence). While Hilal-i-Imtiaz, Pride of Performance, Sitara-i-Imtiaz, and Tamgha-e-Imtiaz occupies a unique role and importance in Pakistan's civil society. Atta-ur-Rahman is the only scientist of Pakistan to have won all these 4 Civil Awards.
Achievements:
In 1961, international achievements first recorded in 1961 when Pakistan became the third Asian country and tenth in the world when the Rehbar-I – a solid fuel expendable rocket— was launched from Sonmani Spaceport. The Rehbar-I was developed and launched under the leadership of Dr. W. J. M. Turowicz, a Polish-Pakistani scientist and then project director of this program. Since then, the program began taking flights which continued until the 1970s.
Achievements:
A major breakthrough occurred in 1979, when the Nobel Prize Committee awarded the Nobel Prize in Physics to Abdus Salam, for formulating the electroweak theory – a theory that provides the basis of unification of weak nuclear force and electromagnetic force. In 1990, the Space and Upper Atmosphere Research Commission (SUPARCO) launched the first, and locally designed, a communication satellite, Badr-1, from Xichang Satellite Launch Center (XLSC) of the People's Republic of China. With the launch, Pakistan became the first Muslim majority country to have developed an artificial robotic satellite, and was the second South Asian state to have launched its satellite, second to India.
Achievements:
One of the widely reported achievements was in 1998 when the country joined the nuclear club. In response to India's nuclear tests on 11 May and 13 May 1998, under codename Operation Shakti, in the long-constructed Pokhran Test Range (PTR). Under Prime Minister Nawaz Sharif, the Pakistan Atomic Energy Commission (PAEC) conducted five simultaneous tests at the Chagai Hills under codename Chagai-I on 28 May 1998. PAEC carried out another test in the Kharan Desert, under Chagai-II, meaning it had tested six devices in under one week. With the testing of these atomic devices, Pakistan became the seventh nuclear power in the world, and the only Muslim-majority country to have mastered the technology. On 13 August 2011, SUPARCO launched its first indigenously developed geosynchronous satellite, Paksat-1R also from XLSC in China.
Achievements:
In 2006 Prof. Atta-ur-Rahman was elected as Fellow of Royal Society (London), thereby becoming the first scientist from the Muslim world to be so honoured in recognition of researches and contributions carried out within an Islamic country. He has major contributions in the development of natural product chemistry and several international journals have published special issue in recognition of these contributions in his honour, He contributed to the major development of science and technology as Chairman Higher Education Commission during 2002–2008 which have resulted in a significant increase in research publications in Pakistan from only about 800 research papers in Impact Factor journals in 2002 to over 11,000 publications in 2016 the quality of which has been recognised by ThomsonReuters. The International Centre for Chemical and Biological Sciences at the University of Karachi which has developed as a leading research centre in the region under the leadership of Prof. Atta-ur-Rahman was designated as a UNESCO Centre of Excellence in 2016. Prof. Atta-ur-Rahman was awarded the high Civil Award of the Government of Austria (the 'Grosses goldenes Ehrenzeichen am Bande') in 2007 in recognition for his contributions for uplifting science in Pakistan, and the Government of China also honoured him with the highest Award for Foreigners (Friendship Award) in recognition of his eminent contributions. The largest university of Malaysia, Universiti Teknologi Mara, established a Research Centre entitled " Dr. Atta-ur-Rahman Research Institute of natural Product Discovery" to honour this great Muslim scientist for uplifting science in Pakistan and in the Muslim world in his capacity as Coordinator General COMSTECH, a Ministerial Committee comprising 57 Ministers of Science and Technology of the 57 OIC member countries. More recently, the leading Chinese University on Traditional Medicine in Changsha, Hunan has also decided to name a research institute in honour of Prof. Atta-ur-Rahman FRS, in recognition of his eminent contributions to uplift science in Pakistan and to establish strong linkages with China.In another landmark study undertaken by Thomson Reuters, highlighting the impact of the reforms introduced by Atta-ur-Rahman, it was revealed that the rate of growth of highly cited papers from Pakistan in a decade was even greater than that in Brazil, Russia, India or ChinaIn reply to C.N.R. Rao Professor Atta-ur-Rahman wrote recently, The Indian government need not be worried. We Pakistanis, alas, know how to destroy our own institutions.
Information technology:
The rapid progress made by Pakistan in the IT and telecom sector during 2000–2002, under Professor Atta-ur-Rahman as Federal Minister, led to the spread of internet from 29 cities in the year 2000 to 1,000 cities, towns and villages by 2002, and the spread of fiber from 40 cities to 400 cities in this period. The first IT policy and implementation strategy was approved under the leadership of Prof. Atta-ur-Rahman, then Federal Minister of Science & technology, in August 2000 which laid the foundations of the development of this sector The internet prices were reduced sharply from $87,000 per month for a 2 MB line to only $3000 per month and later to $90 per month. The mobile telephony boom also occurred under the leadership of Atta-ur-Rahman, and it began by the drastic lowering of prices, bringing in of competition (Ufone) and changing the system so that the person receiving a call was no longer required to pay any charges. A satellite was placed in space (Paksat 1) at a cost of only $4 million. These changes in the IT infra-structure proved invaluable for the Higher education sector. Pakistan Educational Research Network was set up in 2004 through which one of the finest digital libraries was established in universities. In 2002, few university libraries could subscribe to a handful of journals. Today every student in every public sector university has free access to over 20,000 international journals with back volumes and over 60,000 books from 250 international publishers.
Information technology:
As of 2011, Pakistan has over 20 million internet users and is ranked as one of the top countries that have registered a high growth rate in internet penetration. Overall, it has the 15th largest population of internet users in the world. In the fiscal year 2012–2013, the Government of Pakistan aims to spend 4.6 billion rupees (Rs.) on information technology projects, with emphasis on e-government, human resource and infrastructure development.
Information technology:
Pakistan's information technology industry has gone through a dramatic change, and the country has taken the lead in adopting some technologies while also setting an example for others in global best practices. Matters relating to the IT industry are overseen and regulated by the Ministry of Information Technology of the Government of Pakistan. The IT industry is regarded as a successful sector of Pakistan economically, even during the financial crisis. The Government of Pakistan has given numerous favors to IT investors in the country since the last decade, that resulted in the development of the IT sector. In the years 2003–2005 the country's IT exports saw a rise of about fifty percent and amounted a total of about US$48.5 million. The World Economic Forum, assessing the development of Information and Communication Technology in the country ranked Pakistan 102nd among 144 countries in the Global Information Technology report of 2012.
Higher education reforms:
Reform 2002–2009 In 2002, the University Grants Commission was replaced by the Higher Education Commission (HEC), which has an independent chairperson. The HEC was charged with reforming Pakistan's higher education system by introducing better financial incentives, increasing university enrolment and the number of PhD graduates, boosting foreign scholarships and research collaboration and providing all the major universities with state-of-the-art ICT facilities.In a series of reforms in 2002, the HEC instituted major upgrades for scientific laboratories, rehabilitating existing educational facilities, expanding research support and overseeing the development of one of the best digital libraries in the region. Seeking to meet international standards, quality assurance and accreditation process was also established. Some ~95% of students sent abroad for training returned, an unusually high result for a developing country, in response to improved salaries and working conditions at universities as well as bonding and strict follow-up by the commission, Fulbright and others. Within a limited timespan, the HEC provided all universities with free, high-speed Internet access to scientific literature, an upgrade of research equipment accessible across the country and a programme for the creation of new universities of science and technology, including science parks which attracted foreign investors.
Higher education reforms:
International praise : Pakistan's Golden Period for Higher Education Since the Higher Education Commission (HEC) reforms have been carried out in 2002, HEC has received praise from international higher education observers. Rahman, founding Chairman of HEC, has received a number of international awards for the transformation of the higher education sector under his leadership.[25] German academic, Dr. Wolfgang Voelter of Tübingen University in Germany over viewed the performance of HEC under the leadership of Rahman and described the reforms in HEC as "A miracle happened." After teaching and visiting in 15 universities of Pakistan, Voelter wrote that the "scenario of education, science and technology in Pakistan has changed dramatically, as never before in the history of the country." The chairperson of the Senate Standing Committee on Education recently announced the first 6 years of HEC under Rahman as "Pakistan's golden period in higher education".American academic Prof. Fred M. Hayward has also praised the reform process undertaken by Pakistan, admitting that "since 2002, a number of extraordinary changes have taken place." Hayward pointed out that "over the last six years almost 4,000 scholars have participated in Ph.D. programs in Pakistan in which more than 600 students have studied in foreign PhD programs'.The HEC's reforms were also applauded by the United Nations Commission on Science and Technology for Development (UNCSTD) which reported that the "progress made was breath-taking and has put Pakistan ahead of comparable countries in numerous aspects." The UNCSTD has closely monitored the development in Pakistan in the past years, coming to the conclusion that HEC's program initiated under the leadership of Rahman is a "best-practice" example for developing countries aiming at building their human resources and establishing an innovative, technology-based economy.". According to an article published in the leading science journal Nature "Rahman's strong scientific background, enthusiasm for reform and impressive ability to secure cash made him a hit at home and abroad. It really was an anomaly that we had a person of that stature with that kind of backing,----Atta-ur-Rahman was a force of natureRahman has won four international awards for the revolutionary changes in the higher education sector brought in the HEC. Nature, a leading science journal, has also written a number of editorials and articles about the transformation brought about in Pakistan in the higher education sector under the HEC. In an article entitled "Pakistan Threat to Indian Science" published in the leading daily newspaper Hindustan Times, India, it has been reported that Professor C. N. R. Rao, Chairman of the Indian Prime Minister's Scientific Advisory Council made a presentation to the Indian Prime Minister at the rapid progress made by Pakistan in the higher education sector under the leadership of Rahman, Chairman, Higher Education Commission. It was reported that as a result of the reforms, "Pakistan may soon join China in giving India serious competition in science". "Science is a lucrative profession in Pakistan. It has tripled the salaries of its scientists in the last few years." Decentralizing the governance of higher education In 2011–2012, the HEC found itself on the brink of dissolution in the face of the 18th amendment to the Constitution, which devolved several governance functions to provincial governments, including that of higher education. It was only after Prof. Atta-ur-Rahman FRS former Chairman HEC filed a petition before the Supreme Court of Pakistan and the Supreme Court intervened in April 2011, that the commission was spared from being divided up among the four Provinces of Baluchistan, Khyber–Pakhtunkhwa, Punjab and Sindh.Notwithstanding this, the HEC's developmental budget – that spent on scholarships and faculty training, etc. – was slashed by 37.8% in 2011–2012, from a peak of R. 22.5 billion (circa US$0.22 billion) in 2009–2010 to Rs 14 billion (circa US$0.14 billion). The higher education sector continues to face an uncertain future, despite the marginal increase in developmental spending wrought by the new administration in Islamabad: Rs. 18.5 billion (circa US$0.18 billion) in the 2013–2014 budget. According to HEC statistics, the organization's budget as a percentage of national GDP has consistently fallen from the 2006–2007 peak of 0.33% to 0.19% in 2011–2012.In defiance of the Supreme Court ruling of April 2011, the provincial assembly of Sindh Province passed the unprecedented Sindh Higher Commission Act in 2013 creating Pakistan's first provincial higher education commission. In October 2014, Punjab Province followed suit as part of a massive restructuring of its own higher education system.
Higher education reforms:
Effect of reforms on student numbers and academic output Despite the turbulence caused by the legal battle being waged since the 2011 constitutional amendment discussed above, the number of degree-awarding institutions continues to grow throughout the country, both in the private and public sectors. University student rolls have continued to rise, from 0.28 million in 2001 to 0.47 million in 2005 and more than 1.2 million in 2014. Just under half of universities are privately owned.Between 2002 and 2009, the HEC increased the number of PhD graduates to 6 000 per year and in provided up to 11 000 scholarships for study abroad. The number of Pakistani publications recorded in Thomson Reuter's Web of Science (Science Citation Index Expanded) leapt from 714 to 3 614 over the same period then to 6778 by 2014, and to over 20,000 by 2020. This progress in scientific productivity appears to be due to the momentum generated by the larger numbers of faculty and student scholarships for study abroad, as well as the swelling ranks of PhD graduates. Critics argue that the rapid, massive increase in numbers has compromised quality. However this claim has been refuted by neutral international experts.
Challenges:
Pakistan has been known for some of its achievements in science and technology such as successful development of media and military technologies and a growing base of doctors and engineers, as well as its new influx of software engineers who have been contributing to Pakistan's Information Technology industry. Due to present situation in Pakistan, around 3,000 Pakistani doctors emigrate to Western economies in search of suitable employment opportunities and hence contribute intellectually to the health sector of developed countries and at the same time leaving the effects of a brain drain in Pakistan.Pervez Hoodbhoy published a report on scientific output in Pakistan in which he claimed that research and scientific activities are lower than many other developing countries Hoodbhoy asserted that Pakistan has produced fewer papers than neighboring India. The contentions of Hoodbhoy have been questioned for using outdated data. The increase in research output from Pakistan has been praised after the establishment of the Higher Education Commission in 2002. This is borne out by the graphical comparison between Pakistan and India shown on the right which shows that Pakistan (green) was 400% behind India (blue) in research publications per 10 million population in year 2000 but overtook India in 2017 and by 2018, it was about 20% ahead of India according to Web of Science data.
Challenges:
In a report published by Thomson Reuters in 2016, it has been concluded that the rate of increase of highly cited papers in international journals from Pakistan is higher than that from Brazil, Russia, India or China.Pakistan’s public-sector infrastructure for science and technology is complemented by academic institutions and the strategic and defence sectors. Over the years, these three components have vied for political patronage and societal recognition, leading to duplication and competition between the different bodies.
Scientific research institutions (SRI):
A large part of research is conducted by science research institutes with semi-controlled by the Government. International Center for Chemical and Biological Sciences H.E.J. Research Institute of Chemistry Dr. Panjwani Centre for Molecular Medicine and Drug Research School of Biological Sciences, Punjab University National Center for Physics National Institute for Biotechnology and Genetic Engineering Abdus Salam School of Mathematical Sciences PU Centre for High Energy Physics Atta-ur-Rahman School of Applied Biosciences, NUST Institute of Space and Planetary Astrophysics National Engineering and Scientific Commission Pakistan Institute of Nuclear Science and Technology Institute of Space Technology Council of Scientific and Industrial Research Nuclear Institute for Agriculture and Biology Nuclear Institute for Food and Agriculture Technology Resource Mobilization Unit Federal Bureau of Statistics Mathematics Statistical Division
Science community of Pakistan:
NUST Science Society Pakistan Mathematical Society Pakistan Agricultural Research Council Pakistan Academy of Sciences Pakistan Institute of Physics Pakistan Astrophysicist Society Pakistan Atomic Energy Commission Pakistan Atomic Scientists Society Pakistan Nuclear Society National Information and Communication Technologies Research and Development Funds Pakistan Science Foundation Department of Pakistan Survey Pakistan Geo-engineering and Geological Survey Pakistan Cave Research & Caving Federation Pakistan Physical Society Pakistan Optical Society Khwarizmi Science Society Pakistan science club Ghulam Ishaq Khan Institute of Engineering Sciences and Technology Shaheed Zulfiqar Ali Bhutto Institute of Science and Technology Pakistan Institute of Nuclear Science and Technology National Institute of Food Science and Technology USTAD Institute of Science & Technology Abbottabad Royal Institute of Science & Technology Karachi Gandhara Institute of Science & Technology Sukkur Institute of Science & Technology Bright Institute of Science and technology - Peshawar Pakistan Advanced Institute of Science and Technology
Sources:
This article incorporates text from a free content work. Licensed under CC-BY-SA IGO 3.0. Text taken from UNESCO Science Report: towards 2030, UNESCO, UNESCO Publishing. To learn how to add open license text to Wikipedia articles, please see this how-to page. For information on reusing text from Wikipedia, please see the terms of use.
Sources:
This article incorporates text from a free content work. Licensed under CC BY-SA 3.0 IGO (license statement/permission). Text taken from UNESCO Science Report: the Race Against Time for Smarter Development, 574–603, UNESCO Publishing. To learn how to add open license text to Wikipedia articles, please see this how-to page. For information on reusing text from Wikipedia, please see the terms of use. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**CD146**
CD146:
CD146 (cluster of differentiation 146) also known as the melanoma cell adhesion molecule (MCAM) or cell surface glycoprotein MUC18, is a 113kDa cell adhesion molecule currently used as a marker for endothelial cell lineage. In humans, the CD146 protein is encoded by the MCAM gene.
Function:
MCAM functions as a receptor for laminin alpha 4, a matrix molecule that is broadly expressed within the vascular wall. Accordingly, MCAM is highly expressed by cells that are components of the blood vessel wall, including vascular endothelial cells, smooth muscle cells and pericytes. Its function is still poorly understood, but evidence points to it being part of the endothelial junction associated with the actin cytoskeleton. A member of the Immunoglobulin superfamily, it consists of five Ig domains, a transmembrane domain, and a cytoplasmic region. It is expressed on chicken embryonic spleen and thymus, activated human T cells, endothelial progenitors such as angioblasts and mesenchymal stem cells, and strongly expressed on blood vessel endothelium and smooth muscle.
Function:
Two isoforms exist (MCAM long (MCAM-1), and MCAM short, or MCAM-s) which differ in the length of their cytoplasmic domain. Activation of these isoforms seems to produce functional differences as well. Natural killer cells transfected with MCAM-1 demonstrate decreased rolling velocity and increased cell adhesion to an endothelial cell monolayer and increased microvilli formation while cells transfected with MCAM-s showed no change in adhesion characteristics. Since these characteristics are important in leukocyte extravasation, MCAM-1 may be an important part of the inflammatory response.
Function:
CD146 has been demonstrated to appear on a small subset of T and B lymphocytes in the peripheral blood of healthy individuals. The CD146+ T cells display an immunophenotype consistent with effector memory cells and have a distinct gene profile from the CD146- T cells. CD146 T cells have been shown by Dagur and colleagues to produce IL-17.CD146 has been seen as a marker for mesenchymal stem cells isolated from multiple adult and fetal organs, and its expression may be linked to multipotency; mesenchymal stem cells with greater differentiation potential express higher levels of CD146 on the cell surface.
Relevance in cancer:
MCAM inhibits breast cancer progression.Normal melanocytes do not express MCAM and the expression of MCAM is first found in nevi and melanoma cells. MCAM expression is positively correlated to melanoma progression at which the expression of MCAM is highest in metastatic melanoma cells. The significance of MCAM upregulation is evident in melanoma cells cultured in 3D skin reconstruct in which MCAM facilitates the migration of melanoma into the dermis. Without the expression of MCAM melanoma cells are controlled by keratinocytes in the epidermis that inhibit penetrance beyond the basement membrane. The control by keratinocytes are only achieved by E-cadherin expression on the surface of melanoma cells. Melanoma cells with functional E-cadherin on the surface can only exclusively grow in the epidermis as keratinocytes frequently downregulate the expression of MCAM on melanoma cells. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Holmquistite**
Holmquistite:
Holmquistite is a lithium magnesium aluminium inosilicate mineral with chemical formula Li2(Mg,Fe2+)3Al2Si8O22(OH)2. It crystallizes in the orthorhombic crystal system as prismatic crystals up to 10 cm (3.9 in) or as massive aggregates. It has a Mohs hardness of 5-6 and a specific gravity of 2.95 to 3.13.
Color could vary from black, dark violet to light sky blue.
It occurs as metasomatic replacements on the margins of lithium rich pegmatites.
It was first described in 1913 from an occurrence in Utö, near Stockholm, Sweden. It was named for the Swedish petrologist Per Johan Holmquist (1866–1946). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Best of Tests DS**
Best of Tests DS:
Best of Tests DS is an educational video game released for Nintendo DS in North America on March 4, 2008. The game intends to present tests and puzzles that will help stimulate one's Intelligent Quotient or "IQ" through five categories: logic, observation, memory, speed of perception and analysis. The game contains three gaming modes: Training, Memory and Tests. Best of Tests DS calculates the player's IQ and adapts the difficulty level in order to provide a personalized challenge that is always renewed and updated.
Reception:
Best of Tests DS received negative reviews from critics upon release. On Metacritic, the game holds a score of 32/100 based on 8 reviews, indicating "generally unfavorable reviews." | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**ORPHEUS**
ORPHEUS:
The Organisation for PhD Education in Biomedicine and Health Sciences in the European System ("ORPHEUS") is an organization committed to safeguard the PhD as a research degree and to strengthen career opportunities for PhD graduates.
History:
ORPHEUS was established in April 2004 with a European Conference on Harmonisation of PhD Programmes in Biomedicine and Health Sciences held in Zagreb. It became clear that despite many similarities PhD programmes there are also important differences in content and the standard expected. Thus a PhD title (also known as Doctor of Philosophy) may have different meanings in different parts of Europe. While some countries did not have any clinical PhD programmes, others had even two. Thus, the delegates coming from 25 universities and from 16 European countries agreed that there was a need for European harmonisation and they accepted 'The Declaration of the European Conference on Harmonisation of PhD Programmes in Medicine and Health Sciences' known as the 'Zagreb Declaration', which contain the first European consensus statement what a PhD programme should consist of and aim for. Many of the ideas of the 'Zagreb Declaration' are also seen in the Salzburg Principles and subsequent Salzburg II recommendations from the EUA-CDE.
History:
Since 2004, ORPHEUS conferences have been held annually at different countries within Europe (see below), with normal attendances of 150-200 participants.
Aims:
ORPHEUS has the following aims: To give active support and guidance to members of ORPHEUS in enhancing their contributions to medicine and society in general.
To provide information to members of ORPHEUS and all PhD candidates all over Europe.
To represent higher education and research in biomedicine and health sciences and to influence policy making at national, European and international level To encourage cooperation among members of the Association and the development of effective bilateral and multilateral networks.
To promote cooperation in research and development of joint PhD programmes.
To promote harmonisation of PhD programmes in biomedicine and health sciences To encourage mobility of PhD candidates and academic staff.
To stimulate quality assurance of PhD research and education, and in particular to develop an accreditation process of PhD programmes in biomedicine and health sciences To cooperate with other associations with similar goals
ORPHEUS Self-evaluation and Labelling:
The primary ORPHEUS tool is its Best Practices which in a concise manner provides recommendations for the aims and content of PhD programmes in biomedicine and health sciences. The document is the result of extensive consultation throughout Europe with stakeholders including deans, graduate school heads, supervisors and students. The document has been found to be appropriate for institutions throughout Europe.
ORPHEUS Self-evaluation and Labelling:
Implementation of the Best Practices is encouraged through the ORPHEUS Labelling programme that allows institutions to self-evaluate and obtain an ORPHEUS Label. Through a simple questionnaire, institutions can determine the extent to which their PhD programmes are consistent with the ORPHEUS recommendations as shown in the Best Practices document. Institutions which believe they comply with these recommendations can apply for a Label. In this case an evaluation team is set up by the Labelling Board to assess if this case usually in connection with a site-visit. Institutions who do not qualify for a Label may be awarded an Evaluation Certificate.
Key ORPHEUS documents:
Best Practices for PhD Training Standards for PhD education in Biomedicine and Health Sciences in Europe. 2012 Towards Standards for PhD Education in Biomedicine and Health Sciences. A position paper from ORPHEUS. 2009 Helsinki Consensus Statement on PhD Training in Clinical Research. 2007 'Zagreb Declaration. 2004'
ORPHEUS conferences:
ORPHEUS 2020 conference delayed to 2021 due to Covid-19 pandemic.
ORPHEUS conferences:
ORPHEUS 2019, 14th European Conference Dublin, 2019 ORPHEUS 2018, 13th European Conference Reykjavik, 2018 ORPHEUS 2017, 12th European Conference Klaipeda, 2017 ORPHEUS 2016, 11th European Conference Cologne, 2016 ORPHEUS 2015, 10th European Conference Belgrade, 2015 ORPHEUS 2014, Ninth European Conference Lausanne, 2014 ORPHEUS 2013, Eighth European Conference Prague, 2013 ORPHEUS 2012. Seventh European Conference. Establishing Evaluation of PhD Training. Bergen 2012.
ORPHEUS conferences:
ORPHEUS 2011. Sixth European Conference. PhD Quality Indicators for Biomedicine and Health Sciences. Izmir 2011.
ORPHEUS 2010. Fifth European Conference: The Advancement of European Biomedical and Health Science PhD Education by Cooperative Networking. Vienna 2010.
ORPHEUS2009. Fourth European Conference: Setting Standards for PhD Education in Biomedicine and Health Sciences. Aarhus 2009.
ORPHEUS2007. Third European Conference: Biomedical and Health Science Doctoral Training. Helsinki 2007.
Second European Conference on Harmonisation of PhD Programmes in Biomedicine and Health Sciences. Zagreb 2005.
First European Conference on Harmonization of PhD programs in Biomedicine and Health Sciences. Zagreb 2004. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Ikitoxin**
Ikitoxin:
Ikitoxin is a neurotoxin from the venom of the South African Spitting scorpion (Parabuthus transvaalicus) that targets voltage-sensitive sodium channels. It causes unprovoked jumps in mice following intracerebroventricular injections.
Sources:
Ikitoxin is one of the many components that can be isolated from the venom of the South African Spitting scorpion. Other peptide toxins found in the venom include birtoxin, which is moderately toxic but very abundant, dortoxin, a lethal peptide, bestoxin, which causes writhing in mice, and altitoxin, a highly depressant peptide.
Chemistry:
Ikitoxin is a member of the birtoxin family of peptide neurotoxins that target sodium channels. Although identified as a long chain neurotoxin, which usually have 64-70 residues with four disulfide bridges, ikitoxin, like birtoxin, has a smaller size (58 residues) with only three disulfide bridges. Ikitoxin differs from birtoxin by a single amino acid: from glycine to glutamic acid at position 23, consistent with an apparent mass difference of 72 Da between the two peptides.
Mode of action:
Both ikitoxin and birtoxin are beta toxins, which bind to and trap the voltage sensor of the channel at side 4. The binding of ikitoxin lowers the voltage threshold of sodium channels and produce a reduction in the current amplitude. As a result of the change in their activation properties, sodium channels will open at smaller depolarizations, resulting in increased excitability.
Toxicity:
Ikitoxin differs from birtoxin by only a single residue, but has a markedly reduced biological activity. In mice experiments, intracerebroventricular administration of ikitoxin induced unprovoked jumps. These jumps were observed at a concentration that was 1000-fold higher than in the case of birtoxin, and their onset was much slower. Another difference between these toxins is that birtoxin produced convulsions, tremors, increased ventilation and, subsequently, death. Injection of up to 4 μg of ikitoxin in mice was not lethal. Ikitoxin seems to affect only mammals.
Treatment:
Ikitoxin is one of many neurotoxic polypeptide components in the venom of the South African Spitting scorpion. It has a birtoxin-like structure. Antibodies against the N-terminus of the birtoxin protein structure can neutralize the venom of the South African spitting scorpion, and such antibodies may be useful clinically to treat envenomation. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Zinc finger and scan domain containing 30**
Zinc finger and scan domain containing 30:
Zinc finger and SCAN domain containing 30 is a protein that in humans is encoded by the ZSCAN30 gene. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**GJA1**
GJA1:
Gap junction alpha-1 protein (GJA1), also known as connexin 43 (Cx43), is a protein that in humans is encoded by the GJA1 gene on chromosome 6. As a connexin, GJA1 is a component of gap junctions, which allow for gap junction intercellular communication (GJIC) between cells to regulate cell death, proliferation, and differentiation. As a result of its function, GJA1 is implicated in many biological processes, including muscle contraction, embryonic development, inflammation, and spermatogenesis, as well as diseases, including oculodentodigital dysplasia (ODDD), heart malformations, and cancers.
Structure:
GJA1 is a 43.0 kDa protein composed of 382 amino acids. GJA1 contains a long C-terminal tail, an N-terminal domain, and multiple transmembrane domains. The protein passes through the phospholipid bilayer four times, leaving its C- and N-terminals exposed to the cytoplasm. The C-terminal tail is composed of 50 amino acids and includes post-translational modification sites, as well as binding sites for transcription factors, cytoskeleton elements, and other proteins. As a result, the C-terminal tail is central to functions such as regulating pH gating and channel assembly. Notably, the DNA region of the GJA1 gene encoding this tail is highly conserved, indicating that it is either resistant to mutations or becomes lethal when mutated. Meanwhile, the N-terminal domain is involved in channel gating and oligomerization and, thus, may control the switch between the channel's open and closed states. The transmembrane domains form the gap junction channel while the extracellular loops facilitate proper channel docking. Moreover, two extracellular loops form disulfide bonds that interact with two hexamers to form a complete gap junction channel.The connexin-43 internal ribosome entry site is an RNA element present in the 5' UTR of the mRNA of GJA1. This internal ribosome entry site (IRES) allows cap independent translation during conditions such as heat shock and stress.
Function:
As a member of the connexin family, GJA1 is a component of gap junctions, which are intercellular channels that connect adjacent cells to permit the exchange of low molecular weight molecules, such as small ions and secondary messengers, to maintain homeostasis.GJA1 is the most ubiquitously expressed connexin and is detected in most cell types.
Function:
It is the major protein in heart gap junctions and is purported to play a crucial role in the synchronized contraction of the heart. Despite its key role in the heart and other vital organs, GJA1 has a short half-life (only two to four hours), indicating that the protein undergoes daily turnover in the heart and may be highly abundant or compensated with other connexins. GJA1 is also largely involved in embryonic development. For instance, transforming growth factor-beta 1 (TGF-β1) was observed to induce GJA1 expression via the Smad and ERK1/2 signaling pathways, resulting in trophoblast cell differentiation into the placenta.Furthermore, GJA1 is expressed in many immune cells, such as eosinophils and T cells, where its gap junction function promotes the maturation and activation of these cells and, by extension, the cross-communication necessary to mount an inflammatory response. It has also been shown that uterine macrophage directly physically couple with uterine myocytes through GJA1, transferring Ca²⁺, to promote uterine muscle contraction and excitation during human labor onset.In addition, GJA1 can be found in the Leydig cells and seminiferous tubules between Sertoli cells and spermatogonia or primary spermatocytes, where it plays a key role in spermatogenesis and testis development through controlling the tight junction proteins in the blood-testis barrier.
Function:
While it is a channel protein, GJA1 can also perform channel-independent functions. In the cytoplasm, the protein regulates the microtubule network and, by extension, cell migration and polarity. This function has been observed in brain and heart development, as well as wound-healing in endothelial cells. GJA1 has also been observed to localize to the mitochondria, where it promotes cell survival by downregulating the intrinsic apoptotic pathway during conditions of oxidative stress.
Clinical significance:
Mutations in this gene have been associated with ODDD; craniometaphyseal dysplasia; sudden infant death syndrome, which is linked to cardiac arrhythmia; Hallermann–Streiff syndrome; and heart malformations, such as viscero-atrial heterotaxia. There have also been a few cases of reported hearing loss and skin disorders unrelated to ODDD. Ultimately, GJA1 has low tolerance for deviations from its original sequence, with mutations resulting in loss- or gain-of-channel function that lead to disease phenotypes. It is paradoxical, however, that patients with an array of somatic mutations in GJA1 most often do not present with cardiac arrhythmias, even though connexin-43 is the most abundant protein forming gap junctional pores in cardiomyocytes and are essential for normal action potential propagation.Notably, GJA1 expression has been associated with a wide variety of cancers, including nasopharyngeal carcinoma, meningioma, hemangiopericytoma, liver tumor, colon cancer, esophageal cancer, breast cancer, mesothelioma, glioblastoma, lung cancer, adrenocortical tumors, renal cell cancer, cervical carcinoma, ovarian carcinoma, endometrial carcinoma, prostate cancer, thyroid carcinoma, and testicular cancer. Its role in controlling cell motility and polarity was thought to contribute to cancer development and metastasis, though its role as a gap junction protein may also be involved. Moreover, the cytoprotective effects of this protein can promote tumor cell survival in radiotherapy treatments, while silencing its gene increases radiosensitivity. As a result, GJA1 may serve as a target for improving the success of radiotherapeutic treatment of cancer. As a biomarker, GJA1 could also be used to screen young males for risk of testis cancer.The thyroid hormone triiodothyronine (T3) downregulates the expression of GJA1. This is assumed to be a key mechanism why the conduction velocity in myocardial tissue is reduced in thyrotoxicosis, thereby promoting cardiac arrhythmia.Currently, only rotigaptide, an antiarrhythmic peptide-based drug, and its derivatives, such as danegaptide, have reached clinical trials for treating cardiac pathologies by enhancing GJA1 expression. Alternatively, drugs could target complementary connexins, such as Cx40, which function similarly to GJA1. However, both approaches still require a system to target the diseased tissue to avoid inducing developmental abnormalities elsewhere. Thus, a more effective approach entails designing a miRNA through antisense oligonucleotides, transfection, or infection to knock down only mutant GJA1 mRNA, thus allowing the expression of wildtype GJA1 and retaining normal phenotype.
Interactions:
Gap junction protein, alpha 1 has been shown to interact with: Cx37, Cx40, Cx45, MAPK7, Caveolin 1, Tight junction protein 1 CSNK1D, and PTPmu (PTPRM). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Propellant mass fraction**
Propellant mass fraction:
In aerospace engineering, the propellant mass fraction is the portion of a vehicle's mass which does not reach the destination, usually used as a measure of the vehicle's performance. In other words, the propellant mass fraction is the ratio between the propellant mass and the initial mass of the vehicle. In a spacecraft, the destination is usually an orbit, while for aircraft it is their landing location. A higher mass fraction represents less weight in a design. Another related measure is the payload fraction, which is the fraction of initial weight that is payload. It can be applied to a vehicle, a stage of a vehicle or to a rocket propulsion system.
Formulation:
The propellant mass fraction is given by: where: ζ is the propellant mass fraction m0=mf+mp is the initial mass of the vehicle mp is the propellant mass mf is the final mass of the vehicle
Significance:
In rockets for a given target orbit, a rocket's mass fraction is the portion of the rocket's pre-launch mass (fully fueled) that does not reach orbit. The propellant mass fraction is the ratio of just the propellant to the entire mass of the vehicle at takeoff (propellant plus dry mass). In the cases of a single-stage-to-orbit (SSTO) vehicle or suborbital vehicle, the mass fraction equals the propellant mass fraction, which is simply the fuel mass divided by the mass of the full spaceship. A rocket employing staging, which are the only designs to have reached orbit, has a mass fraction higher than the propellant mass fraction because parts of the rocket itself are dropped off en route. Propellant mass fractions are typically around 0.8 to 0.9.
Significance:
In aircraft, mass fraction is related to range, an aircraft with a higher mass fraction can go farther. Aircraft mass fractions are typically around 0.5.
Significance:
When applied to a rocket as a whole, a low mass fraction is desirable, since it indicates a greater capability for the rocket to deliver payload to orbit for a given amount of fuel. Conversely, when applied to a single stage, where the propellant mass fraction calculation doesn't include the payload, a higher propellant mass fraction corresponds to a more efficient design, since there is less non-propellant mass. Without the benefit of staging, SSTO designs are typically designed for mass fractions around 0.9. Staging increases the payload fraction, which is one of the reasons SSTOs appear difficult to build.
Significance:
For example, the complete Space Shuttle system has: fueled weight at liftoff: 1,708,500 kg dry weight at liftoff: 342,100 kgGiven these numbers, the propellant mass fraction is 342 100 kg 708 500 kg 0.7998 The mass fraction plays an important role in the rocket equation: ln mfm0 Where mf/m0 is the ratio of final mass to initial mass (i.e., one minus the mass fraction), Δv is the change in the vehicle's velocity as a result of the fuel burn and ve is the effective exhaust velocity (see below).
Significance:
The term effective exhaust velocity is defined as: sp where Isp is the fuel's specific impulse in seconds and gn is the standard acceleration of gravity (note that this is not the local acceleration of gravity).
To make a powered landing from orbit on a celestial body without an atmosphere requires the same mass reduction as reaching orbit from its surface, if the speed at which the surface is reached is zero. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Multi-antimicrobial extrusion protein**
Multi-antimicrobial extrusion protein:
Multi-antimicrobial extrusion protein (MATE) also known as multidrug and toxin extrusion or multidrug and toxic compound extrusion is a family of proteins which function as drug/sodium or proton antiporters.
Function:
The MATE proteins in bacteria, archaea and eukaryotes function as fundamental transporters of metabolic and xenobiotic organic cations.
Structure:
These proteins are predicted to have 12 alpha-helical transmembrane regions, some of the animal proteins may have an additional C-terminal helix. The X-ray structure of the NorM was determined to 3.65 Å, revealing an outward-facing conformation with two portals open to the outer leaflet of the membrane and a unique topology of the predicted 12 transmembrane helices distinct from any other known multidrug resistance transporter.
Discovery:
The multidrug efflux transporter NorM from V. parahaemolyticus which mediates resistance to multiple antimicrobial agents (norfloxacin, kanamycin, ethidium bromide etc.) and its homologue from E. coli were identified in 1998. NorM seems to function as drug/sodium antiporter which is the first example of Na+-coupled multidrug efflux transporter discovered. NorM is a prototype of a new transporter family and Brown et al. named it the multidrug and toxic compound extrusion family. NorM is nicknamed "Last of the multidrug transporters" because it is the last multidrug transporter discovered functionally as well as structurally.
Genes:
The following human genes encode MATE proteins: SLC47A1 SLC47A2 | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Tragic triad**
Tragic triad:
The tragic triad is a term used in logotherapy, coined by Dr. Viktor Frankl. The tragic triad refers to three experiences which often lead to existential crisis, namely, guilt, suffering or death. The concept of the tragic triad is used in identifying the life meanings of patients, or the relatives of patients, experiencing guilt, suffering or death. These life meanings are analyzed using logotherapy's existential analysis with the intent of assisting the patient overcome their existential crisis by discovering meaning or purpose in the experience.Frankl argued that all human beings at one point in their lives will encounter the tragic triad. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**HSPA7**
HSPA7:
Heat shock 70kDa protein 7 (HSP70B) also known as HSPA7 is a human gene. The protein encoded by this gene is a member of the Hsp70 family of heat shock proteins. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**3,4-Methylenedioxyphentermine**
3,4-Methylenedioxyphentermine:
3,4-Methylenedioxyphentermine (MDPH) is a lesser-known psychedelic drug. MDPH was first synthesized by Alexander Shulgin. In his book PiHKAL (Phenethylamines i Have Known And Loved), the dosage range is listed as 160–240 mg, and the duration as 3–5 hours. MDPH's effects are very similar to those of MDA: they both are smooth and "stoning," and do not cause any visuals. They also alter dreams and dream patterns. Shulgin describes MDPH as a promoter; it promotes the effects of other drugs, similarly to 2C-D. Very little data exists about the pharmacological properties, metabolism, and toxicity of MDPH.
Legality:
United Kingdom This substance is a Class A drug in the Drugs controlled by the UK Misuse of Drugs Act. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Rocket engine**
Rocket engine:
A rocket engine uses stored rocket propellants as the reaction mass for forming a high-speed propulsive jet of fluid, usually high-temperature gas. Rocket engines are reaction engines, producing thrust by ejecting mass rearward, in accordance with Newton's third law. Most rocket engines use the combustion of reactive chemicals to supply the necessary energy, but non-combusting forms such as cold gas thrusters and nuclear thermal rockets also exist. Vehicles propelled by rocket engines are commonly called rockets. Rocket vehicles carry their own oxidiser, unlike most combustion engines, so rocket engines can be used in a vacuum to propel spacecraft and ballistic missiles.
Rocket engine:
Compared to other types of jet engine, rocket engines are the lightest and have the highest thrust, but are the least propellant-efficient (they have the lowest specific impulse). The ideal exhaust is hydrogen, the lightest of all elements, but chemical rockets produce a mix of heavier species, reducing the exhaust velocity.
Rocket engines become more efficient at high speeds, due to the Oberth effect.
Terminology:
Here, "rocket" is used as an abbreviation for "rocket engine".
Thermal rockets use an inert propellant, heated by electricity (electrothermal propulsion) or a nuclear reactor (nuclear thermal rocket).
Chemical rockets are powered by exothermic reduction-oxidation chemical reactions of the propellant: Solid-fuel rockets (or solid-propellant rockets or motors) are chemical rockets which use propellant in a solid state.
Liquid-propellant rockets use one or more propellants in a liquid state fed from tanks.
Hybrid rockets use a solid propellant in the combustion chamber, to which a second liquid or gas oxidiser or propellant is added to permit combustion.
Monopropellant rockets use a single propellant decomposed by a catalyst. The most common monopropellants are hydrazine and hydrogen peroxide.
Principle of operation:
Rocket engines produce thrust by the expulsion of an exhaust fluid that has been accelerated to high speed through a propelling nozzle. The fluid is usually a gas created by high pressure (150-to-4,350-pound-per-square-inch (10 to 300 bar)) combustion of solid or liquid propellants, consisting of fuel and oxidiser components, within a combustion chamber. As the gases expand through the nozzle, they are accelerated to very high (supersonic) speed, and the reaction to this pushes the engine in the opposite direction. Combustion is most frequently used for practical rockets, as the laws of thermodynamics (specifically Carnot's theorem) dictate that high temperatures and pressures are desirable for the best thermal efficiency. Nuclear thermal rockets are capable of higher efficiencies, but currently have environmental problems which preclude their routine use in the Earth's atmosphere and cislunar space.
Principle of operation:
For model rocketry, an available alternative to combustion is the water rocket pressurized by compressed air, carbon dioxide, nitrogen, or any other readily available, inert gas.
Propellant Rocket propellant is mass that is stored, usually in some form of tank, or within the combustion chamber itself, prior to being ejected from a rocket engine in the form of a fluid jet to produce thrust.
Chemical rocket propellants are the most commonly used. These undergo exothermic chemical reactions producing a hot gas jet for propulsion. Alternatively, a chemically inert reaction mass can be heated by a high-energy power source through a heat exchanger in lieu of a combustion chamber.
Solid rocket propellants are prepared in a mixture of fuel and oxidising components called grain, and the propellant storage casing effectively becomes the combustion chamber.
Principle of operation:
Injection Liquid-fuelled rockets force separate fuel and oxidiser components into the combustion chamber, where they mix and burn. Hybrid rocket engines use a combination of solid and liquid or gaseous propellants. Both liquid and hybrid rockets use injectors to introduce the propellant into the chamber. These are often an array of simple jets – holes through which the propellant escapes under pressure; but sometimes may be more complex spray nozzles. When two or more propellants are injected, the jets usually deliberately cause the propellants to collide as this breaks up the flow into smaller droplets that burn more easily.
Principle of operation:
Combustion chamber For chemical rockets the combustion chamber is typically cylindrical, and flame holders, used to hold a part of the combustion in a slower-flowing portion of the combustion chamber, are not needed. The dimensions of the cylinder are such that the propellant is able to combust thoroughly; different rocket propellants require different combustion chamber sizes for this to occur.
Principle of operation:
This leads to a number called L∗ , the characteristic length: L∗=VcAt where: Vc is the volume of the chamber At is the area of the throat of the nozzle.L* is typically in the range of 64–152 centimetres (25–60 in).
Principle of operation:
The temperatures and pressures typically reached in a rocket combustion chamber in order to achieve practical thermal efficiency are extreme compared to a non-afterburning airbreathing jet engine. No atmospheric nitrogen is present to dilute and cool the combustion, so the propellant mixture can reach true stoichiometric ratios. This, in combination with the high pressures, means that the rate of heat conduction through the walls is very high.
Principle of operation:
In order for fuel and oxidiser to flow into the chamber, the pressure of the propellants entering the combustion chamber must exceed the pressure inside the combustion chamber itself. This may be accomplished by a variety of design approaches including turbopumps or, in simpler engines, via sufficient tank pressure to advance fluid flow. Tank pressure may be maintained by several means, including a high-pressure helium pressurization system common to many large rocket engines or, in some newer rocket systems, by a bleed-off of high-pressure gas from the engine cycle to autogenously pressurize the propellant tanks For example, the self-pressurization gas system of the SpaceX Starship is a critical part of SpaceX strategy to reduce launch vehicle fluids from five in their legacy Falcon 9 vehicle family to just two in Starship, eliminating not only the helium tank pressurant but all hypergolic propellants as well as nitrogen for cold-gas reaction-control thrusters.
Principle of operation:
Nozzle The hot gas produced in the combustion chamber is permitted to escape through an opening (the "throat"), and then through a diverging expansion section. When sufficient pressure is provided to the nozzle (about 2.5–3 times ambient pressure), the nozzle chokes and a supersonic jet is formed, dramatically accelerating the gas, converting most of the thermal energy into kinetic energy. Exhaust speeds vary, depending on the expansion ratio the nozzle is designed for, but exhaust speeds as high as ten times the speed of sound in air at sea level are not uncommon. About half of the rocket engine's thrust comes from the unbalanced pressures inside the combustion chamber, and the rest comes from the pressures acting against the inside of the nozzle (see diagram). As the gas expands (adiabatically) the pressure against the nozzle's walls forces the rocket engine in one direction while accelerating the gas in the other.
Principle of operation:
The most commonly used nozzle is the de Laval nozzle, a fixed geometry nozzle with a high expansion-ratio. The large bell- or cone-shaped nozzle extension beyond the throat gives the rocket engine its characteristic shape.
Principle of operation:
The exit static pressure of the exhaust jet depends on the chamber pressure and the ratio of exit to throat area of the nozzle. As exit pressure varies from the ambient (atmospheric) pressure, a choked nozzle is said to be under-expanded (exit pressure greater than ambient), perfectly expanded (exit pressure equals ambient), over-expanded (exit pressure less than ambient; shock diamonds form outside the nozzle), or grossly over-expanded (a shock wave forms inside the nozzle extension).In practice, perfect expansion is only achievable with a variable-exit area nozzle (since ambient pressure decreases as altitude increases), and is not possible above a certain altitude as ambient pressure approaches zero. If the nozzle is not perfectly expanded, then loss of efficiency occurs. Grossly over-expanded nozzles lose less efficiency, but can cause mechanical problems with the nozzle. Fixed-area nozzles become progressively more under-expanded as they gain altitude. Almost all de Laval nozzles will be momentarily grossly over-expanded during startup in an atmosphere.Nozzle efficiency is affected by operation in the atmosphere because atmospheric pressure changes with altitude; but due to the supersonic speeds of the gas exiting from a rocket engine, the pressure of the jet may be either below or above ambient, and equilibrium between the two is not reached at all altitudes (see diagram).
Principle of operation:
Back pressure and optimal expansion For optimal performance, the pressure of the gas at the end of the nozzle should just equal the ambient pressure: if the exhaust's pressure is lower than the ambient pressure, then the vehicle will be slowed by the difference in pressure between the top of the engine and the exit; on the other hand, if the exhaust's pressure is higher, then exhaust pressure that could have been converted into thrust is not converted, and energy is wasted.
Principle of operation:
To maintain this ideal of equality between the exhaust's exit pressure and the ambient pressure, the diameter of the nozzle would need to increase with altitude, giving the pressure a longer nozzle to act on (and reducing the exit pressure and temperature). This increase is difficult to arrange in a lightweight fashion, although is routinely done with other forms of jet engines. In rocketry a lightweight compromise nozzle is generally used and some reduction in atmospheric performance occurs when used at other than the 'design altitude' or when throttled. To improve on this, various exotic nozzle designs such as the plug nozzle, stepped nozzles, the expanding nozzle and the aerospike have been proposed, each providing some way to adapt to changing ambient air pressure and each allowing the gas to expand further against the nozzle, giving extra thrust at higher altitudes.
Principle of operation:
When exhausting into a sufficiently low ambient pressure (vacuum) several issues arise. One is the sheer weight of the nozzle—beyond a certain point, for a particular vehicle, the extra weight of the nozzle outweighs any performance gained. Secondly, as the exhaust gases adiabatically expand within the nozzle they cool, and eventually some of the chemicals can freeze, producing 'snow' within the jet. This causes instabilities in the jet and must be avoided.
Principle of operation:
On a de Laval nozzle, exhaust gas flow detachment will occur in a grossly over-expanded nozzle. As the detachment point will not be uniform around the axis of the engine, a side force may be imparted to the engine. This side force may change over time and result in control problems with the launch vehicle.
Advanced altitude-compensating designs, such as the aerospike or plug nozzle, attempt to minimize performance losses by adjusting to varying expansion ratio caused by changing altitude.
Principle of operation:
Propellant efficiency For a rocket engine to be propellant efficient, it is important that the maximum pressures possible be created on the walls of the chamber and nozzle by a specific amount of propellant; as this is the source of the thrust. This can be achieved by all of: heating the propellant to as high a temperature as possible (using a high energy fuel, containing hydrogen and carbon and sometimes metals such as aluminium, or even using nuclear energy) using a low specific density gas (as hydrogen rich as possible) using propellants which are, or decompose to, simple molecules with few degrees of freedom to maximise translational velocitySince all of these things minimise the mass of the propellant used, and since pressure is proportional to the mass of propellant present to be accelerated as it pushes on the engine, and since from Newton's third law the pressure that acts on the engine also reciprocally acts on the propellant, it turns out that for any given engine, the speed that the propellant leaves the chamber is unaffected by the chamber pressure (although the thrust is proportional). However, speed is significantly affected by all three of the above factors and the exhaust speed is an excellent measure of the engine propellant efficiency. This is termed exhaust velocity, and after allowance is made for factors that can reduce it, the effective exhaust velocity is one of the most important parameters of a rocket engine (although weight, cost, ease of manufacture etc. are usually also very important).
Principle of operation:
For aerodynamic reasons the flow goes sonic ("chokes") at the narrowest part of the nozzle, the 'throat'. Since the speed of sound in gases increases with the square root of temperature, the use of hot exhaust gas greatly improves performance. By comparison, at room temperature the speed of sound in air is about 340 m/s while the speed of sound in the hot gas of a rocket engine can be over 1700 m/s; much of this performance is due to the higher temperature, but additionally rocket propellants are chosen to be of low molecular mass, and this also gives a higher velocity compared to air.
Principle of operation:
Expansion in the rocket nozzle then further multiplies the speed, typically between 1.5 and 2 times, giving a highly collimated hypersonic exhaust jet. The speed increase of a rocket nozzle is mostly determined by its area expansion ratio—the ratio of the area of the exit to the area of the throat, but detailed properties of the gas are also important. Larger ratio nozzles are more massive but are able to extract more heat from the combustion gases, increasing the exhaust velocity.
Principle of operation:
Thrust vectoring Vehicles typically require the overall thrust to change direction over the length of the burn. A number of different ways to achieve this have been flown: The entire engine is mounted on a hinge or gimbal and any propellant feeds reach the engine via low pressure flexible pipes or rotary couplings.
Just the combustion chamber and nozzle is gimballed, the pumps are fixed, and high pressure feeds attach to the engine.
Multiple engines (often canted at slight angles) are deployed but throttled to give the overall vector that is required, giving only a very small penalty.
High-temperature vanes protrude into the exhaust and can be tilted to deflect the jet.
Overall performance:
Rocket technology can combine very high thrust (meganewtons), very high exhaust speeds (around 10 times the speed of sound in air at sea level) and very high thrust/weight ratios (>100) simultaneously as well as being able to operate outside the atmosphere, and while permitting the use of low pressure and hence lightweight tanks and structure.
Rockets can be further optimised to even more extreme performance along one or more of these axes at the expense of the others.
Overall performance:
Specific impulse The most important metric for the efficiency of a rocket engine is impulse per unit of propellant, this is called specific impulse (usually written Isp ). This is either measured as a speed (the effective exhaust velocity ve in metres/second or ft/s) or as a time (seconds). For example, if an engine producing 100 pounds of thrust runs for 320 seconds and burns 100 pounds of propellant, then the specific impulse is 320 seconds. The higher the specific impulse, the less propellant is required to provide the desired impulse.
Overall performance:
The specific impulse that can be achieved is primarily a function of the propellant mix (and ultimately would limit the specific impulse), but practical limits on chamber pressures and the nozzle expansion ratios reduce the performance that can be achieved.
Net thrust Below is an approximate equation for calculating the net thrust of a rocket engine: Since, unlike a jet engine, a conventional rocket motor lacks an air intake, there is no 'ram drag' to deduct from the gross thrust. Consequently, the net thrust of a rocket motor is equal to the gross thrust (apart from static back pressure).
Overall performance:
The m˙ve−opt term represents the momentum thrust, which remains constant at a given throttle setting, whereas the Ae(pe−pamb) term represents the pressure thrust term. At full throttle, the net thrust of a rocket motor improves slightly with increasing altitude, because as atmospheric pressure decreases with altitude, the pressure thrust term increases. At the surface of the Earth the pressure thrust may be reduced by up to 30%, depending on the engine design. This reduction drops roughly exponentially to zero with increasing altitude.
Overall performance:
Maximum efficiency for a rocket engine is achieved by maximising the momentum contribution of the equation without incurring penalties from over expanding the exhaust. This occurs when pe=pamb . Since ambient pressure changes with altitude, most rocket engines spend very little time operating at peak efficiency.
Since specific impulse is force divided by the rate of mass flow, this equation means that the specific impulse varies with altitude.
Overall performance:
Vacuum specific impulse, Isp Due to the specific impulse varying with pressure, a quantity that is easy to compare and calculate with is useful. Because rockets choke at the throat, and because the supersonic exhaust prevents external pressure influences travelling upstream, it turns out that the pressure at the exit is ideally exactly proportional to the propellant flow m˙ , provided the mixture ratios and combustion efficiencies are maintained. It is thus quite usual to rearrange the above equation slightly: and so define the vacuum Isp to be: where: And hence: Throttling Rockets can be throttled by controlling the propellant combustion rate m˙ (usually measured in kg/s or lb/s). In liquid and hybrid rockets, the propellant flow entering the chamber is controlled using valves, in solid rockets it is controlled by changing the area of propellant that is burning and this can be designed into the propellant grain (and hence cannot be controlled in real-time).
Overall performance:
Rockets can usually be throttled down to an exit pressure of about one-third of ambient pressure (often limited by flow separation in nozzles) and up to a maximum limit determined only by the mechanical strength of the engine.
In practice, the degree to which rockets can be throttled varies greatly, but most rockets can be throttled by a factor of 2 without great difficulty; the typical limitation is combustion stability, as for example, injectors need a minimum pressure to avoid triggering damaging oscillations (chugging or combustion instabilities); but injectors can be optimised and tested for wider ranges.
For example, some more recent liquid-propellant engine designs that have been optimised for greater throttling capability (BE-3, Raptor) can be throttled to as low as 18–20 per cent of rated thrust.
Solid rockets can be throttled by using shaped grains that will vary their surface area over the course of the burn.
Overall performance:
Energy efficiency Rocket engine nozzles are surprisingly efficient heat engines for generating a high speed jet, as a consequence of the high combustion temperature and high compression ratio. Rocket nozzles give an excellent approximation to adiabatic expansion which is a reversible process, and hence they give efficiencies which are very close to that of the Carnot cycle. Given the temperatures reached, over 60% efficiency can be achieved with chemical rockets.
Overall performance:
For a vehicle employing a rocket engine the energetic efficiency is very good if the vehicle speed approaches or somewhat exceeds the exhaust velocity (relative to launch); but at low speeds the energy efficiency goes to 0% at zero speed (as with all jet propulsion). See Rocket energy efficiency for more details.
Thrust-to-weight ratio Rockets, of all the jet engines, indeed of essentially all engines, have the highest thrust-to-weight ratio. This is especially true for liquid-fuelled rocket engines.
Overall performance:
This high performance is due to the small volume of pressure vessels that make up the engine—the pumps, pipes and combustion chambers involved. The lack of inlet duct and the use of dense liquid propellant allows the pressurisation system to be small and lightweight, whereas duct engines have to deal with air which has around three orders of magnitude lower density.
Overall performance:
Of the liquid fuels used, density is lowest for liquid hydrogen. Although hydrogen/oxygen burning has the highest specific impulse of any in-use chemical rocket, hydrogen's very low density (about one-fourteenth that of water) requires larger and heavier turbopumps and pipework, which decreases the engine's thrust-to-weight ratio (for example the RS-25) compared to those that do not use hydrogen (NK-33).
Mechanical issues:
Rocket combustion chambers are normally operated at fairly high pressure, typically 10–200 bar (1–20 MPa, 150–3,000 psi). When operated within significant atmospheric pressure, higher combustion chamber pressures give better performance by permitting a larger and more efficient nozzle to be fitted without it being grossly overexpanded.
However, these high pressures cause the outermost part of the chamber to be under very large hoop stresses – rocket engines are pressure vessels.
Worse, due to the high temperatures created in rocket engines the materials used tend to have a significantly lowered working tensile strength.
In addition, significant temperature gradients are set up in the walls of the chamber and nozzle, these cause differential expansion of the inner liner that create internal stresses.
Acoustic issues:
The extreme vibration and acoustic environment inside a rocket motor commonly result in peak stresses well above mean values, especially in the presence of organ pipe-like resonances and gas turbulence.
Acoustic issues:
Combustion instabilities The combustion may display undesired instabilities, of sudden or periodic nature. The pressure in the injection chamber may increase until the propellant flow through the injector plate decreases; a moment later the pressure drops and the flow increases, injecting more propellant in the combustion chamber which burns a moment later, and again increases the chamber pressure, repeating the cycle. This may lead to high-amplitude pressure oscillations, often in ultrasonic range, which may damage the motor. Oscillations of ±200 psi at 25 kHz were the cause of failures of early versions of the Titan II missile second stage engines. The other failure mode is a deflagration to detonation transition; the supersonic pressure wave formed in the combustion chamber may destroy the engine.Combustion instability was also a problem during Atlas development. The Rocketdyne engines used in the Atlas family were found to suffer from this effect in several static firing tests, and three missile launches exploded on the pad due to rough combustion in the booster engines. In most cases, it occurred while attempting to start the engines with a "dry start" method whereby the igniter mechanism would be activated prior to propellant injection. During the process of man-rating Atlas for Project Mercury, solving combustion instability was a high priority, and the final two Mercury flights sported an upgraded propulsion system with baffled injectors and a hypergolic igniter.
Acoustic issues:
The problem affecting Atlas vehicles was mainly the so-called "racetrack" phenomenon, where burning propellant would swirl around in a circle at faster and faster speeds, eventually producing vibration strong enough to rupture the engine, leading to complete destruction of the rocket. It was eventually solved by adding several baffles around the injector face to break up swirling propellant.
More significantly, combustion instability was a problem with the Saturn F-1 engines. Some of the early units tested exploded during static firing, which led to the addition of injector baffles.
Acoustic issues:
In the Soviet space program, combustion instability also proved a problem on some rocket engines, including the RD-107 engine used in the R-7 family and the RD-216 used in the R-14 family, and several failures of these vehicles occurred before the problem was solved. Soviet engineering and manufacturing processes never satisfactorily resolved combustion instability in larger RP-1/LOX engines, so the RD-171 engine used to power the Zenit family still used four smaller thrust chambers fed by a common engine mechanism.
Acoustic issues:
The combustion instabilities can be provoked by remains of cleaning solvents in the engine (e.g. the first attempted launch of a Titan II in 1962), reflected shock wave, initial instability after ignition, explosion near the nozzle that reflects into the combustion chamber, and many more factors. In stable engine designs the oscillations are quickly suppressed; in unstable designs they persist for prolonged periods. Oscillation suppressors are commonly used.
Acoustic issues:
Three different types of combustion instabilities occur: Chugging A low frequency oscillation in chamber pressure below 200 Hertz. Usually it is caused by pressure variations in feed lines due to variations in acceleration of the vehicle, when rocket engines are building up thrust, are shut down or are being throttled.: 261 : 146 Chugging can cause a worsening feedback loop, as cyclic variation in thrust causes longitudinal vibrations to travel up the rocket, causing the fuel lines to vibrate, which in turn do not deliver propellant smoothly into the engines. This phenomenon is known as "pogo oscillations" or "pogo", named after the pogo stick.: 258 In the worst case, this may result in damage to the payload or vehicle. Chugging can be minimised by using several methods, such as installing energy-absorbing devices on feed lines.: 259 Chugging may cause Screeching.: 146 Buzzing An intermediate frequency oscillation in chamber pressure between 200 and 1000 Hertz. Usually caused due to insufficient pressure drop across the injectors.: 261 It generally is mostly annoying, rather than being damaging.
Acoustic issues:
Buzzing is known to have adverse effects on engine performance and reliability, primarily as it causes material fatigue.: 147 In extreme cases combustion can end up being forced backwards through the injectors – this can cause explosions with monopropellants. Buzzing may cause Screeching.: 261 Screeching A high frequency oscillation in chamber pressure above 1000 Hertz, sometimes called screaming or squealing. The most immediately damaging, and the hardest to control. It is due to acoustics within the combustion chamber that often couples to the chemical combustion processes that are the primary drivers of the energy release, and can lead to unstable resonant "screeching" that commonly leads to catastrophic failure due to thinning of the insulating thermal boundary layer. Acoustic oscillations can be excited by thermal processes, such as the flow of hot air through a pipe or combustion in a chamber. Specifically, standing acoustic waves inside a chamber can be intensified if combustion occurs more intensely in regions where the pressure of the acoustic wave is maximal. Such effects are very difficult to predict analytically during the design process, and have usually been addressed by expensive, time-consuming and extensive testing, combined with trial and error remedial correction measures.
Acoustic issues:
Screeching is often dealt with by detailed changes to injectors, changes in the propellant chemistry, vaporising the propellant before injection or use of Helmholtz dampers within the combustion chambers to change the resonant modes of the chamber.Testing for the possibility of screeching is sometimes done by exploding small explosive charges outside the combustion chamber with a tube set tangentially to the combustion chamber near the injectors to determine the engine's impulse response and then evaluating the time response of the chamber pressure- a fast recovery indicates a stable system.
Acoustic issues:
Exhaust noise For all but the very smallest sizes, rocket exhaust compared to other engines is generally very noisy. As the hypersonic exhaust mixes with the ambient air, shock waves are formed. The Space Shuttle generated over 200 dB(A) of noise around its base. To reduce this, and the risk of payload damage or injury to the crew atop the stack, the mobile launcher platform was fitted with a Sound Suppression System that sprayed 1.1 million litres (290,000 US gal) of water around the base of the rocket in 41 seconds at launch time. Using this system kept sound levels within the payload bay to 142 dB.The sound intensity from the shock waves generated depends on the size of the rocket and on the exhaust velocity. Such shock waves seem to account for the characteristic crackling and popping sounds produced by large rocket engines when heard live. These noise peaks typically overload microphones and audio electronics, and so are generally weakened or entirely absent in recorded or broadcast audio reproductions. For large rockets at close range, the acoustic effects could actually kill.More worryingly for space agencies, such sound levels can also damage the launch structure, or worse, be reflected back at the comparatively delicate rocket above. This is why so much water is typically used at launches. The water spray changes the acoustic qualities of the air and reduces or deflects the sound energy away from the rocket.
Acoustic issues:
Generally speaking, noise is most intense when a rocket is close to the ground, since the noise from the engines radiates up away from the jet, as well as reflecting off the ground. Also, when the vehicle is moving slowly, little of the chemical energy input to the engine can go into increasing the kinetic energy of the rocket (since useful power P transmitted to the vehicle is P=F∗V for thrust F and speed V). Then the largest portion of the energy is dissipated in the exhaust's interaction with the ambient air, producing noise. This noise can be reduced somewhat by flame trenches with roofs, by water injection around the jet and by deflecting the jet at an angle.
Rocket engine development:
United States The development of the US rocket engine industry has been shaped by a complex web of relationships between government agencies, private companies, research institutions, and other stakeholders. Since the establishment of the first liquid-propellant rocket engine company (Reaction Motors, Inc.) in 1941 and the first government laboratory (GALCIT) devoted to the subject, the US liquid-propellant rocket engine (LPRE) industry has undergone significant changes. At least 14 US companies have been involved in the design, development, manufacture, testing, and flight support operations of various types of rocket engines from 1940 to 2000. In contrast to other countries like Russia, China, or India, where only government or pseudogovernment organisations engage in this business, the US government relies heavily on private industry. These commercial companies are essential to the continued viability of the United States and its form of governance, as they compete with one another to provide cutting-edge rocket engines that meet the needs of the government, the military, and the private sector. In the United States the company that develops the LPRE usually is awarded the production contract. Generally, the need or demand for a new rocket engine comes from government agencies such as NASA or the Department of Defense. Once the need is identified, government agencies may issue requests for proposals (RFPs) to solicit proposals from private companies and research institutions. Private companies and research institutions, in turn, may invest in research and development (R&D) activities to develop new rocket engine technologies that meet the needs and specifications outlined in the RFPs.
Rocket engine development:
Alongside private companies, universities, independent research institutes and government laboratories also play a critical role in the research and development of rocket engines. Universities provide graduate and undergraduate education to train qualified technical personnel, and their research programs often contribute to the advancement of rocket engine technologies. More than 25 universities in the US have taught or are currently teaching courses related to Liquid Propellant Rocket Engines (LPREs), and their graduate and undergraduate education programs are considered one of their most important contributions. Universities such as Princeton University, Cornell University, Purdue University, Pennsylvania State University, University of Alabama, the Navy's Post-Graduate School, or the California Institute of Technology have conducted excellent R&D work on topics related to the rocket engine industry. One of the earliest examples of the contribution of universities to the rocket engine industry is the work of the GALCIT in 1941. They demonstrated the first jet-assisted takeoff (JATO) rockets to the Army, leading to the establishment of the Jet Propulsion Laboratory.
Rocket engine development:
However the transfer of knowledge from research professors and their projects to the rocket engine industry has been a mixed experience. While some notable professors and relevant research projects have positively influenced industry practices and understanding of LPREs, the connection between university research and commercial companies has been inconsistent and weak. Universities were not always aware of the industry's specific needs, and engineers and designers in the industry had limited knowledge of university research. As a result, many university research programs remained relatively unknown to industry decision-makers. Furthermore, in the last few decades, certain university research projects, while interesting to professors, were not useful to the industry due to a lack of communication or relevance to industry needs.Government laboratories, including the Rocket Propulsion Laboratory (now part of Air Force Research Laboratory), Arnold Engineering Test Center, NASA Marshall Space Flight Center, Jet Propulsion Laboratory, Stennis Space Center, White Sands Proving Grounds, and NASA John H. Glenn Research Center, have played crucial roles in the development of liquid rocket propulsion engines (LPREs). They have conducted unbiased testing, guided work at US and some non-US contractors, performed research and development, and provided essential testing facilities including hover test facilities and simulated altitude test facilities and resources. Initially, private companies or foundations financed smaller test facilities, but since the 1950s, the U.S. government has funded larger test facilities at government laboratories. This approach reduced costs for the government by not building similar facilities at contractors' plants but increased complexity and expenses for contractors. Nonetheless, government laboratories have solidified their significance and contributed to LPRE advancements.
Rocket engine development:
LPRE programs have been subject to several cancellations in the United States, even after spending millions of dollars on their development. For example, the M-l LOX/LH2 LPRE, Titan I, and the RS-2200 aerospike, as well as several JATO units and large uncooled thrust chambers were cancelled. The cancellations of these programs were not related to the specific LPRE's performance or any issues with it. Instead, they were due to the cancellation of the vehicle programs the engine was intended for or budget cuts imposed by the government.
Rocket engine development:
USSR Russia and the former Soviet Union was and still is the world's foremost nation in developing and building rocket engines. From 1950 to 1998, their organisations developed, built, and put into operation a larger number and a larger variety of liquid propellant rocket engine (LPRE) designs than any other country. Approximately 500 different LPREs have been developed before 2003. For comparison the United States has developed slightly more than 300 (before 2003). The Soviets also had the most rocket-propelled flight vehicles. They had more liquid propellant ballistic missiles and more space launch vehicles derived or converted from these decommissioned ballistic missiles than any other nation. As of the end of 1998, the Russians (or earlier the Soviet Union) had successfully launched 2573 satellites with LPREs or almost 65% of the world total of 3973. All of these vehicle flights were made possible by the timely development of suitable high-performance reliable LPREs.
Rocket engine development:
Institutions and Actors Unlike many other countries where the development and production of rocket engines were consolidated within a single organisation, the Soviet Union took a different approach, they established numerous specialised design bureaus (DB) which would compete for development contracts. These design bureaus, or "konstruktorskoye buro" (KB) in Russian were state run organisations which were primarily responsible for carrying out research, development and prototyping of advanced technologies usually related to military hardware, such as turbojet engines, aircraft components, missiles, or space launch vehicles.
Rocket engine development:
Design Bureaus which specialised in rocket engines often possessed the necessary personnel, facilities, and equipment to conduct laboratory tests, flow tests, and ground testing of experimental rocket engines. Some even had specialised facilities for testing very large engines, conducting static firings of engines installed in vehicle stages, or simulating altitude conditions during engine tests. In certain cases, engine testing, certification and quality control were outsourced to other organisations and locations with more suitable test facilities. Many DBs also had housing complexes, gymnasiums, and medical facilities intended to support the needs of their employees and their families.
Rocket engine development:
The Soviet Union's LPRE development effort saw significant growth during the 1960s and reached its peak in the 1970s. This era coincided with the Cold War between the Soviet Union and the United States, characterised by intense competition in spaceflight achievements. Between 14 and 17 Design Bureaus and research institutes were actively involved in developing LPREs during this period. These organisations received relatively steady support and funding due to high military and spaceflight priorities, which facilitated the continuous development of new engine concepts and manufacturing methods.
Rocket engine development:
Once a mission with a new vehicle (missile or spacecraft) was established it was passed on to a design bureau whose role was to oversee the development of the entire rocket. If none of the previously developed rocket engines met the needs of the mission, a new rocket engine with specific requirements would be contracted to another DB specialised in LPRE development (oftentimes each DB had expertise in specific types of LPREs with different applications, propellants, or engine sizes). This meant that the development or design study of a rocket engine was always aimed at a specific application which entailed set requirements.
Rocket engine development:
When it comes to which DBs were awarded contracts for the development of new rocket engines either a single design bureau would be chosen or several design bureaus would be given the same contract which sometimes led to fierce competition between DBs.
Rocket engine development:
When only one DB was picked for the development, it was often the result of the relationship between a vehicle or system's chief designer and the chief designer of a rocket engine specialised DB. If the vehicle’s chief designer was happy with previous work done by a certain design bureau it was not unusual to see continued reliance on that LPRE bureau for that class of engines. For example, all but one of the LPREs for submarine-launched missiles were developed by the same design bureau for the same vehicle development prime contractor. However, when two parallel engine development programs were supported in order to select the superior one for a specific application, several qualified rocket engine models were never used. This luxury of choice was not commonly available in other nations. However, the use of design bureaus also led to certain issues, including program cancellations and duplication. Some major programs were cancelled, resulting in the disposal or storage of previously developed engines. One notable example of duplication and cancellation was the development of engines for the R-9A ballistic missile. Two sets of engines were supported, but ultimately only one set was selected, leaving several perfectly functional engines unused. Similarly, for the ambitious heavy N-l space launch vehicle intended for lunar and planetary missions, the Soviet Union developed and put into production at least two engines for each of the six stages. Additionally, they developed alternate engines for a more advanced N-l vehicle. However, the program faced multiple flight failures, and with the United States' successful moon landing, the program was ultimately cancelled, leaving the Soviet Union with a surplus of newly qualified engines without a clear purpose.
Rocket engine development:
These examples demonstrate the complex dynamics and challenges faced by the Soviet Union in managing the development and production of rocket engines through Design Bureaus.
Rocket engine development:
Accidents The development of rocket engines in the Soviet Union was marked by significant achievements, but it also carried ethical considerations due to numerous accidents and fatalities. From a Science and Technology Studies point of view, the ethical implications of these incidents shed light on the complex relationship between technology, human factors, and the prioritisation of scientific advancement over safety.
Rocket engine development:
The Soviet Union encountered a series of tragic accidents and mishaps in the development and operation of rocket engines. Notably, the USSR holds the unfortunate distinction of having experienced more injuries and deaths resulting from liquid propellant rocket engine (LPRE) accidents than any other country. These incidents brought into question the ethical considerations surrounding the development, testing, and operational use of rocket engines.
Rocket engine development:
One of the most notable disasters occurred in 1960 when the R-16 ballistic missile suffered a catastrophic accident on the launchpad at the Tyuratam launch facility. This incident resulted in the deaths of 124 engineers and military personnel, including Marshal M.I. Nedelin, a former deputy minister of defence. The explosion occurred after the second-stage rocket engine suddenly ignited, causing the fully loaded missile to disintegrate. The explosion resulted from the ignition and explosion of the mixed hypergolic propellants, consisting of nitric acid with additives and UDMH (unsymmetrical dimethylhydrazine).
Rocket engine development:
While the immediate cause of the 1960 accident was attributed to a lack of protective circuits in the missile control unit, the ethical considerations surrounding LPRE accidents in the USSR extend beyond specific technical failures. The secrecy surrounding these accidents, which remained undisclosed for approximately three decades, raises concerns about transparency, accountability, and the protection of human life.
Rocket engine development:
The decision to keep fatal LPRE accidents hidden from the public eye reflects a broader ethical dilemma. The Soviet government, driven by the pursuit of scientific and technological superiority during the Cold War, sought to maintain an image of invincibility and conceal the failures that accompanied their advancements. This prioritisation of national prestige over the well-being and safety of workers raises questions about the ethical responsibility of the state and the organisations involved.
Testing:
Rocket engines are usually statically tested at a test facility before being put into production. For high altitude engines, either a shorter nozzle must be used, or the rocket must be tested in a large vacuum chamber.
Safety:
Rocket vehicles have a reputation for unreliability and danger; especially catastrophic failures. Contrary to this reputation, carefully designed rockets can be made arbitrarily reliable. In military use, rockets are not unreliable. However, one of the main non-military uses of rockets is for orbital launch. In this application, the premium has typically been placed on minimum weight, and it is difficult to achieve high reliability and low weight simultaneously. In addition, if the number of flights launched is low, there is a very high chance of a design, operations or manufacturing error causing destruction of the vehicle.
Safety:
Saturn family (1961–1975) The Rocketdyne H-1 engine, used in a cluster of eight in the first stage of the Saturn I and Saturn IB launch vehicles, had no catastrophic failures in 152 engine-flights. The Pratt and Whitney RL10 engine, used in a cluster of six in the Saturn I second stage, had no catastrophic failures in 36 engine-flights. The Rocketdyne F-1 engine, used in a cluster of five in the first stage of the Saturn V, had no failures in 65 engine-flights. The Rocketdyne J-2 engine, used in a cluster of five in the Saturn V second stage, and singly in the Saturn IB second stage and Saturn V third stage, had no catastrophic failures in 86 engine-flights.
Safety:
Space Shuttle (1981–2011) The Space Shuttle Solid Rocket Booster, used in pairs, caused one notable catastrophic failure in 270 engine-flights.
The RS-25, used in a cluster of three, flew in 46 refurbished engine units. These made a total of 405 engine-flights with no catastrophic in-flight failures. A single in-flight RS-25 engine failure occurred during Space Shuttle Challenger's STS-51-F mission. This failure had no effect on mission objectives or duration.
Cooling:
For efficiency reasons, higher temperatures are desirable, but materials lose their strength if the temperature becomes too high. Rockets run with combustion temperatures that can reach 6,000 °F (3,300 °C; 3,600 K).: 98 Most other jet engines have gas turbines in the hot exhaust. Due to their larger surface area, they are harder to cool and hence there is a need to run the combustion processes at much lower temperatures, losing efficiency. In addition, duct engines use air as an oxidant, which contains 78% largely unreactive nitrogen, which dilutes the reaction and lowers the temperatures. Rockets have none of these inherent combustion temperature limiters.
Cooling:
The temperatures reached by combustion in rocket engines often substantially exceed the melting points of the nozzle and combustion chamber materials (about 1,200 K for copper). Most construction materials will also combust if exposed to high temperature oxidiser, which leads to a number of design challenges. The nozzle and combustion chamber walls must not be allowed to combust, melt, or vaporize (sometimes facetiously termed an "engine-rich exhaust").
Cooling:
Rockets that use common construction materials such as aluminium, steel, nickel or copper alloys must employ cooling systems to limit the temperatures that engine structures experience. Regenerative cooling, where the propellant is passed through tubes around the combustion chamber or nozzle, and other techniques, such as film cooling, are employed to give longer nozzle and chamber life. These techniques ensure that a gaseous thermal boundary layer touching the material is kept below the temperature which would cause the material to catastrophically fail.
Cooling:
Material exceptions that can sustain rocket combustion temperatures to a certain degree are carbon–carbon materials and rhenium, although both are subject to oxidation under certain conditions. Other refractory alloys, such as alumina, molybdenum, tantalum or tungsten have been tried, but were given up on due to various issues.
Materials technology, combined with the engine design, is a limiting factor in chemical rockets.
Cooling:
In rockets, the heat fluxes that can pass through the wall are among the highest in engineering; fluxes are generally in the range of 0.8–80 MW/m2 (0.5-50 BTU/in2-sec).: 98 The strongest heat fluxes are found at the throat, which often sees twice that found in the associated chamber and nozzle. This is due to the combination of high speeds (which gives a very thin boundary layer), and although lower than the chamber, the high temperatures seen there. (See § Nozzle above for temperatures in nozzle).
Cooling:
In rockets the coolant methods include:: 98–99 Ablative: The combustion chamber inside walls are lined with a material that traps heat and carries it away with the exhaust as it vaporizes.
Radiative cooling: The engine is made of one or several refractory materials, which take heat flux until its outer thrust chamber wall glows red- or white-hot, radiating the heat away.
Dump cooling: A cryogenic propellant, usually hydrogen, is passed around the nozzle and dumped. This cooling method has various issues, such as wasting propellant. It is only used rarely.
Regenerative cooling: The fuel (and possibly, the oxidiser) of a liquid rocket engine is routed around the nozzle before being injected into the combustion chamber or preburner. This is the most widely applied method of rocket engine cooling.
Cooling:
Film cooling: The engine is designed with rows of multiple orifices lining the inside wall through which additional propellant is injected, cooling the chamber wall as it evaporates. This method is often used in cases where the heat fluxes are especially high, likely in combination with regenerative cooling. A more efficient subtype of film cooling is transpiration cooling, in which propellant passes through a porous inner combustion chamber wall and transpirates. So far, this method has not seen usage due to various issues with this concept.Rocket engines may also use several cooling methods. Examples: Regeneratively and film cooled combustion chamber and nozzle: V-2 Rocket Engine Regeneratively cooled combustion chamber with a film cooled nozzle extension: Rocketdyne F-1 Engine Regeneratively cooled combustion chamber with an ablatively cooled nozzle extension: The LR-91 rocket engine Ablatively and film cooled combustion chamber with a radiatively cooled nozzle extension: Lunar module descent engine (LMDE), Service propulsion system engine (SPS) Radiatively and film cooled combustion chamber with a radiatively cooled nozzle extension: Deep space storable propellant thrustersIn all cases, another effect that aids in cooling the rocket engine chamber wall is a thin layer of combustion gases (a boundary layer) that is notably cooler than the combustion temperature. Disruption of the boundary layer may occur during cooling failures or combustion instabilities, and wall failure typically occurs soon after.
Cooling:
With regenerative cooling a second boundary layer is found in the coolant channels around the chamber. This boundary layer thickness needs to be as small as possible, since the boundary layer acts as an insulator between the wall and the coolant. This may be achieved by making the coolant velocity in the channels as high as possible.: 105–106 Liquid-fuelled engines are often run fuel-rich, which lowers combustion temperatures. This reduces heat loads on the engine and allows lower cost materials and a simplified cooling system. This can also increase performance by lowering the average molecular weight of the exhaust and increasing the efficiency with which combustion heat is converted to kinetic exhaust energy.
Chemistry:
Rocket propellants require a high energy per unit mass (specific energy), which must be balanced against the tendency of highly energetic propellants to spontaneously explode. Assuming that the chemical potential energy of the propellants can be safely stored, the combustion process results in a great deal of heat being released. A significant fraction of this heat is transferred to kinetic energy in the engine nozzle, propelling the rocket forward in combination with the mass of combustion products released.
Chemistry:
Ideally all the reaction energy appears as kinetic energy of the exhaust gases, as exhaust velocity is the single most important performance parameter of an engine. However, real exhaust species are molecules, which typically have translation, vibrational, and rotational modes with which to dissipate energy. Of these, only translation can do useful work to the vehicle, and while energy does transfer between modes this process occurs on a timescale far in excess of the time required for the exhaust to leave the nozzle.
Chemistry:
The more chemical bonds an exhaust molecule has, the more rotational and vibrational modes it will have. Consequently, it is generally desirable for the exhaust species to be as simple as possible, with a diatomic molecule composed of light, abundant atoms such as H2 being ideal in practical terms. However, in the case of a chemical rocket, hydrogen is a reactant and reducing agent, not a product. An oxidizing agent, most typically oxygen or an oxygen-rich species, must be introduced into the combustion process, adding mass and chemical bonds to the exhaust species.
Chemistry:
An additional advantage of light molecules is that they may be accelerated to high velocity at temperatures that can be contained by currently available materials - the high gas temperatures in rocket engines pose serious problems for the engineering of survivable motors.
Chemistry:
Liquid hydrogen (LH2) and oxygen (LOX, or LO2), are the most effective propellants in terms of exhaust velocity that have been widely used to date, though a few exotic combinations involving boron or liquid ozone are potentially somewhat better in theory if various practical problems could be solved.When computing the specific reaction energy of a given propellant combination, the entire mass of the propellants (both fuel and oxidiser) must be included. The exception is in the case of air-breathing engines, which use atmospheric oxygen and consequently have to carry less mass for a given energy output. Fuels for car or turbojet engines have a much better effective energy output per unit mass of propellant that must be carried, but are similar per unit mass of fuel.
Chemistry:
Computer programs that predict the performance of propellants in rocket engines are available.
Ignition:
With liquid and hybrid rockets, immediate ignition of the propellants as they first enter the combustion chamber is essential.
With liquid propellants (but not gaseous), failure to ignite within milliseconds usually causes too much liquid propellant to be inside the chamber, and if/when ignition occurs the amount of hot gas created can exceed the maximum design pressure of the chamber, causing a catastrophic failure of the pressure vessel.
Ignition:
This is sometimes called a hard start or a rapid unscheduled disassembly (RUD).Ignition can be achieved by a number of different methods; a pyrotechnic charge can be used, a plasma torch can be used, or electric spark ignition may be employed. Some fuel/oxidiser combinations ignite on contact (hypergolic), and non-hypergolic fuels can be "chemically ignited" by priming the fuel lines with hypergolic propellants (popular in Russian engines).
Ignition:
Gaseous propellants generally will not cause hard starts, with rockets the total injector area is less than the throat thus the chamber pressure tends to ambient prior to ignition and high pressures cannot form even if the entire chamber is full of flammable gas at ignition.
Ignition:
Solid propellants are usually ignited with one-shot pyrotechnic devices and combustion usually proceeds through total consumption of the propellants.Once ignited, rocket chambers are self-sustaining and igniters are not needed and combustion usually proceeds through total consumption of the propellants. Indeed, chambers often spontaneously reignite if they are restarted after being shut down for a few seconds. Unless designed for re-ignition, when cooled, many rockets cannot be restarted without at least minor maintenance, such as replacement of the pyrotechnic igniter or even refueling of the propellants.
Jet physics:
Rocket jets vary depending on the rocket engine, design altitude, altitude, thrust and other factors.
Jet physics:
Carbon-rich exhausts from kerosene-based fuels such as RP-1 are often orange in colour due to the black-body radiation of the unburnt particles, in addition to the blue Swan bands. Peroxide oxidiser-based rockets and hydrogen rocket jets contain largely steam and are nearly invisible to the naked eye but shine brightly in the ultraviolet and infrared ranges. Jets from solid-propellant rockets can be highly visible, as the propellant frequently contains metals such as elemental aluminium which burns with an orange-white flame and adds energy to the combustion process. Rocket engines which burn liquid hydrogen and oxygen will exhibit a nearly transparent exhaust, due to it being mostly superheated steam (water vapour), plus some unburned hydrogen.
Jet physics:
The nozzle is usually over-expanded at sea level, and the exhaust can exhibit visible shock diamonds through a schlieren effect caused by the incandescence of the exhaust gas.
The shape of the jet varies for a fixed-area nozzle as the expansion ratio varies with altitude: at high altitude all rockets are grossly under-expanded, and a quite small percentage of exhaust gases actually end up expanding forwards.
Types of rocket engines:
Physically powered Chemically powered Electrically powered Thermal Preheated Solar thermal The solar thermal rocket would make use of solar power to directly heat reaction mass, and therefore does not require an electrical generator as most other forms of solar-powered propulsion do. A solar thermal rocket only has to carry the means of capturing solar energy, such as concentrators and mirrors. The heated propellant is fed through a conventional rocket nozzle to produce thrust. The engine thrust is directly related to the surface area of the solar collector and to the local intensity of the solar radiation and inversely proportional to the Isp.
Types of rocket engines:
Beamed thermal Nuclear thermal Nuclear Nuclear propulsion includes a wide variety of propulsion methods that use some form of nuclear reaction as their primary power source. Various types of nuclear propulsion have been proposed, and some of them tested, for spacecraft applications:
History of rocket engines:
According to the writings of the Roman Aulus Gellius, the earliest known example of jet propulsion was in c. 400 BC, when a Greek Pythagorean named Archytas, propelled a wooden bird along wires using steam. However, it was not powerful enough to take off under its own thrust.
The aeolipile described in the first century BC, often known as Hero's engine, consisted of a pair of steam rocket nozzles mounted on a bearing. It was created almost two millennia before the Industrial Revolution but the principles behind it were not well understood, and it was not developed into a practical power source.
The availability of black powder to propel projectiles was a precursor to the development of the first solid rocket. Ninth Century Chinese Taoist alchemists discovered black powder in a search for the elixir of life; this accidental discovery led to fire arrows which were the first rocket engines to leave the ground.
History of rocket engines:
It is stated that "the reactive forces of incendiaries were probably not applied to the propulsion of projectiles prior to the 13th century". A turning point in rocket technology emerged with a short manuscript entitled Liber Ignium ad Comburendos Hostes (abbreviated as The Book of Fires). The manuscript is composed of recipes for creating incendiary weapons from the mid-eighth to the end of the thirteenth centuries—two of which are rockets. The first recipe calls for one part of colophonium and sulfur added to six parts of saltpeter (potassium nitrate) dissolved in laurel oil, then inserted into hollow wood and lit to "fly away suddenly to whatever place you wish and burn up everything". The second recipe combines one pound of sulfur, two pounds of charcoal, and six pounds of saltpeter—all finely powdered on a marble slab. This powder mixture is packed firmly into a long and narrow case. The introduction of saltpeter into pyrotechnic mixtures connected the shift from hurled Greek fire into self-propelled rocketry. .Articles and books on the subject of rocketry appeared increasingly from the fifteenth through seventeenth centuries. In the sixteenth century, German military engineer Conrad Haas (1509–1576) wrote a manuscript which introduced the construction of multi-staged rockets.Rocket engines were also put in use by Tippu Sultan, the king of Mysore. These usually consisted of a tube of soft hammered iron about 8 in (20 cm) long and 1+1⁄2–3 in (3.8–7.6 cm) diameter, closed at one end, packed with black powder propellant and strapped to a shaft of bamboo about 4 ft (120 cm) long. A rocket carrying about one pound of powder could travel almost 1,000 yards (910 m). These 'rockets', fitted with swords, would travel several meters in the air before coming down with sword edges facing the enemy. These were used very effectively against the British empire.
History of rocket engines:
Modern rocketry Slow development of this technology continued up to the later 19th century, when Russian Konstantin Tsiolkovsky first wrote about liquid-fuelled rocket engines. He was the first to develop the Tsiolkovsky rocket equation, though it was not published widely for some years.
History of rocket engines:
The modern solid- and liquid-fuelled engines became realities early in the 20th century, thanks to the American physicist Robert Goddard. Goddard was the first to use a De Laval nozzle on a solid-propellant (gunpowder) rocket engine, doubling the thrust and increasing the efficiency by a factor of about twenty-five. This was the birth of the modern rocket engine. He calculated from his independently derived rocket equation that a reasonably sized rocket, using solid fuel, could place a one-pound payload on the Moon.
History of rocket engines:
Fritz von Opel was instrumental in popularizing rockets as means of propulsion. In the 1920s, he initiated together with Max Valier, co-founder of the "Verein für Raumschiffahrt", the world's first rocket program, Opel-RAK, leading to speed records for automobiles, rail vehicles and the first manned rocket-powered flight in September of 1929. Months earlier in 1928, one of his rocket-powered prototypes, the Opel RAK2, reached piloted by von Opel himself at the AVUS speedway in Berlin a record speed of 238 km/h, watched by 3000 spectators and world media. A world record for rail vehicles was reached with RAK3 and a top speed of 256 km/h. After these successes, von Opel piloted the world's first public rocket-powered flight using Opel RAK.1, a rocket plane designed by Julius Hatry. World media reported on these efforts, including UNIVERSAL Newsreel of the US, causing as "Raketen-Rummel" or "Rocket Rumble" immense global public excitement, and in particular in Germany, where inter alia Wernher von Braun was highly influenced. The Great Depression led to an end of the Opel-RAK program, but Max Valier continued the efforts. After switching from solid-fuel to liquid-fuel rockets, he died while testing and is considered the first fatality of the dawning space age.
History of rocket engines:
The era of liquid-fuel rocket engines Goddard began to use liquid propellants in 1921, and in 1926 became the first to launch a liquid-fuelled rocket. Goddard pioneered the use of the De Laval nozzle, lightweight propellant tanks, small light turbopumps, thrust vectoring, the smoothly-throttled liquid fuel engine, regenerative cooling, and curtain cooling.: 247–266 During the late 1930s, German scientists, such as Wernher von Braun and Hellmuth Walter, investigated installing liquid-fuelled rockets in military aircraft (Heinkel He 112, He 111, He 176 and Messerschmitt Me 163).The turbopump was employed by German scientists in World War II. Until then cooling the nozzle had been problematic, and the A4 ballistic missile used dilute alcohol for the fuel, which reduced the combustion temperature sufficiently.
History of rocket engines:
Staged combustion (Замкнутая схема) was first proposed by Alexey Isaev in 1949. The first staged combustion engine was the S1.5400 used in the Soviet planetary rocket, designed by Melnikov, a former assistant to Isaev. About the same time (1959), Nikolai Kuznetsov began work on the closed cycle engine NK-9 for Korolev's orbital ICBM, GR-1. Kuznetsov later evolved that design into the NK-15 and NK-33 engines for the unsuccessful Lunar N1 rocket.
History of rocket engines:
In the West, the first laboratory staged-combustion test engine was built in Germany in 1963, by Ludwig Boelkow.
Hydrogen peroxide / kerosene fuelled engines such as the British Gamma of the 1950s used a closed-cycle process by catalytically decomposing the peroxide to drive turbines before combustion with the kerosene in the combustion chamber proper. This gave the efficiency advantages of staged combustion, without the major engineering problems.
History of rocket engines:
Liquid hydrogen engines were first successfully developed in America: the RL-10 engine first flew in 1962. Its successor, the Rocketdyne J-2, was used in the Apollo program's Saturn V rocket to send humans to the Moon. The high specific impulse and low density of liquid hydrogen lowered the upper stage mass and the overall size and cost of the vehicle.
History of rocket engines:
The record for most engines on one rocket flight is 44, set by NASA in 2016 on a Black Brant. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Malvin**
Malvin:
Malvin is a naturally occurring chemical of the anthocyanin family.
Malvin reacts in the presence of H2O2 to form malvone. The ortho-benzoyloxyphenylacetic acid esters reaction product is dependant of the pH: it is obtained under acidic conditions whereas under neutral conditions, the reaction product is the 3-O-acyl-glucosyl-5-O-glucosyl-7-hydroxy coumarin.
Natural occurrences:
It is a diglucoside of malvidin mainly found as a pigment in herbs like Malva (Malva sylvestris), Primula and Rhododendron. M. sylvestris also contains malonylmalvin (malvidin 3-(6″-malonylglucoside)-5-glucoside).The characteristic floral jade coloration of Strongylodon macrobotrys has been shown to be an example of copigmentation, a result of the presence of malvin and saponarin (a flavone glucoside) in the ratio 1:9.
Presence in food Malvin can be found in a variety of common foods, including peaches (Clingstone variety). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**SYT7**
SYT7:
Synaptotagmin-7 is a protein that in humans is encoded by the SYT7 gene.
Function:
Synaptotagmins, such as SYT7, are calcium-dependent phospholipid-binding proteins known for their role in synaptic exocytosis and neurotransmitter release. Significant expression has also been observed in the prostate and other tissues. See MIM 600782 [supplied by OMIM]
Interactions:
SYT7 has been shown to interact with SYNCRIP. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**NDUFA12**
NDUFA12:
NADH dehydrogenase [ubiquinone] 1 alpha subcomplex subunit 12 is an enzyme that in humans is encoded by the NDUFA12 gene. The NDUFA12 protein is a subunit of NADH dehydrogenase (ubiquinone), which is located in the mitochondrial inner membrane and is the largest of the five complexes of the electron transport chain. Mutations in subunits of NADH dehydrogenase (ubiquinone), also known as Complex I, frequently lead to complex neurodegenerative diseases such as Leigh's syndrome that result from mitochondrial complex I deficiency.
Structure:
The NDUFA12 gene is located on the q arm of chromosome 12 in position 22 and spans 32,386 base pairs. The gene produces a 17 kDa protein composed of 145 amino acids. NDUFA12 is a subunit of the enzyme NADH dehydrogenase (ubiquinone), the largest of the respiratory complexes. The structure is L-shaped with a long, hydrophobic transmembrane domain and a hydrophilic domain for the peripheral arm that includes all the known redox centers and the NADH binding site. It has been noted that the N-terminal hydrophobic domain has the potential to be folded into an alpha helix spanning the inner mitochondrial membrane with a C-terminal hydrophilic domain interacting with globular subunits of Complex I. The highly conserved two-domain structure suggests that this feature is critical for the protein function and that the hydrophobic domain acts as an anchor for the NADH dehydrogenase (ubiquinone) complex at the inner mitochondrial membrane. NDUFA12 is one of about 31 hydrophobic subunits that form the transmembrane region of Complex I, but it is an accessory subunit that is believed not to be involved in catalysis. The predicted secondary structure is primarily alpha helix, but the carboxy-terminal half of the protein has high potential to adopt a coiled-coil form. The amino-terminal part contains a putative beta sheet rich in hydrophobic amino acids that may serve as mitochondrial import signal.
Function:
The human NDUFA12 gene codes for a subunit of Complex I of the respiratory chain, which transfers electrons from NADH to ubiquinone. NADH binds to Complex I and transfers two electrons to the isoalloxazine ring of the flavin mononucleotide (FMN) prosthetic arm to form FMNH2. The electrons are transferred through a series of iron-sulfur (Fe-S) clusters in the prosthetic arm and finally to coenzyme Q10 (CoQ), which is reduced to ubiquinol (CoQH2). The flow of electrons changes the redox state of the protein, resulting in a conformational change and pK shift of the ionizable side chain, which pumps four hydrogen ions out of the mitochondrial matrix.
Clinical significance:
Mutations to NDUFA12 are not frequently found to cause complex I deficiency on their own. NDUFA12 is an accessory subunit and the complex can still be found assembled and enzymatically active in its absence, though in reduced amounts and activity. However, a cytosine to tyrosine mutation at position 178 that leads to a premature stop codon has been found in place of arginine at amino acid 60, leading to delayed early development, loss of motor abilities, and basal ganglia lesions typical of Leigh's syndrome. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Silicon controlled rectifier**
Silicon controlled rectifier:
A silicon controlled rectifier or semiconductor controlled rectifier is a four-layer solid-state current-controlling device. The name "silicon controlled rectifier" is General Electric's trade name for a type of thyristor. The principle of four-layer p–n–p–n switching was developed by Moll, Tanenbaum, Goldey, and Holonyak of Bell Laboratories in 1956. The practical demonstration of silicon controlled switching and detailed theoretical behavior of a device in agreement with the experimental results was presented by Dr Ian M. Mackintosh of Bell Laboratories in January 1958. The SCR was developed by a team of power engineers led by Gordon Hall and commercialized by Frank W. "Bill" Gutzwiller in 1957.
Silicon controlled rectifier:
Some sources define silicon-controlled rectifiers and thyristors as synonymous while other sources define silicon-controlled rectifiers as a proper subset of the set of thyristors; the latter being devices with at least four layers of alternating n- and p-type material. According to Bill Gutzwiller, the terms "SCR" and "controlled rectifier" were earlier, and "thyristor" was applied later, as usage of the device spread internationally.SCRs are unidirectional devices (i.e. can conduct current only in one direction) as opposed to TRIACs, which are bidirectional (i.e. charge carriers can flow through them in either direction). SCRs can be triggered normally only by a positive current going into the gate as opposed to TRIACs, which can be triggered normally by either a positive or a negative current applied to its gate electrode.
Modes of operation:
There are three modes of operation for an SCR depending upon the biasing given to it: Forward blocking mode (off state) Forward conduction mode (on state) Reverse blocking mode (off state) Forward blocking mode In this mode of operation, the anode (+, p-doped side) is given a positive voltage while the cathode (−, n-doped side) is given a negative voltage, keeping the gate at zero (0) potential i.e. disconnected. In this case junction J1and J3 are forward-biased, while J2 is reverse-biased, allowing only a small leakage current from the anode to the cathode. When the applied voltage reaches the breakover value for J2, then J2 undergoes avalanche breakdown. At this breakover voltage J2 starts conducting, but below breakover voltage J2 offers very high resistance to the current and the SCR is said to be in the off state.
Modes of operation:
Forward conduction mode An SCR can be brought from blocking mode to conduction mode in two ways: Either by increasing the voltage between anode and cathode beyond the breakover voltage, or by applying a positive pulse at the gate. Once the SCR starts conducting, no more gate voltage is required to maintain it in the ON state. The minimum current necessary to maintain the SCR in the ON state on removal of the gate voltage is called the latching current.
Modes of operation:
There are two ways to turn it off: Reduce the current through it below a minimum value called the holding current, or With the gate turned off, short-circuit the anode and cathode momentarily with a push-button switch or transistor across the junction.
Modes of operation:
Reverse blocking mode When a negative voltage is applied to the anode and a positive voltage to the cathode, the SCR is in reverse blocking mode, making J1 and J3 reverse biased and J2 forward biased. The device behaves as two diodes connected in series. A small leakage current flows. This is the reverse blocking mode. If the reverse voltage is increased, then at critical breakdown level, called the reverse breakdown voltage (VBR), an avalanche occurs at J1 and J3 and the reverse current increases rapidly.
Modes of operation:
SCRs are available with reverse blocking capability, which adds to the forward voltage drop because of the need to have a long, low-doped P1 region. Usually, the reverse blocking voltage rating and forward blocking voltage rating are the same. The typical application for a reverse blocking SCR is in current-source inverters.
Modes of operation:
An SCR incapable of blocking reverse voltage is known as an asymmetrical SCR, abbreviated ASCR. It typically has a reverse breakdown rating in the tens of volts. ASCRs are used where either a reverse conducting diode is applied in parallel (for example, in voltage-source inverters) or where reverse voltage would never occur (for example, in switching power supplies or DC traction choppers).
Modes of operation:
Asymmetrical SCRs can be fabricated with a reverse conducting diode in the same package. These are known as RCTs, for reverse conducting thyristors.
Thyristor turn-on methods:
forward-voltage triggering gate triggering dv/dt triggering thermal triggering light triggeringForward-voltage triggering occurs when the anode–cathode forward voltage is increased with the gate circuit opened. This is known as avalanche breakdown, during which junction J2 will break down. At sufficient voltages, the thyristor changes to its on state with low voltage drop and large forward current. In this case, J1 and J3 are already forward-biased.
Thyristor turn-on methods:
In order for gate triggering to occur, the thyristor should be in the forward blocking state where the applied voltage is less than the breakdown voltage, otherwise forward-voltage triggering may occur. A single small positive voltage pulse can then be applied between the gate and the cathode. This supplies a single gate current pulse that turns the thyristor onto its on state. In practice, this is the most common method used to trigger a thyristor.
Thyristor turn-on methods:
Temperature triggering occurs when the width of depletion region decreases as the temperature is increased. When the SCR is near VPO a very small increase in temperature causes junction J2 to be removed which triggers the device.
Simple SCR circuit:
A simple SCR circuit can be illustrated using an AC voltage source connected to a SCR with a resistive load. Without an applied current pulse to the gate of the SCR, the SCR is left in its forward blocking state. This makes the start of conduction of the SCR controllable. The delay angle α, which is the instant the gate current pulse is applied with respect to the instant of natural conduction (ωt = 0), controls the start of conduction. Once the SCR conducts, the SCR does not turn off until the current through the SCR, is, becomes negative. is stays zero until another gate current pulse is applied and SCR once again begins conducting.
Applications:
SCRs are mainly used in devices where the control of high power, possibly coupled with high voltage, is demanded. Their operation makes them suitable for use in medium- to high-voltage AC power control applications, such as lamp dimming, power regulators and motor control.
Applications:
SCRs and similar devices are used for rectification of high-power AC in high-voltage direct current power transmission. They are also used in the control of welding machines, mainly gas tungsten arc welding and similar processes. It is used as an electronic switch in various devices. Early solid-state pinball machines made use of these to control lights, solenoids, and other functions digitally, instead of mechanically, hence the name solid-state.
Applications:
Other applications include power switching circuits, controlled rectifiers, speed control of DC shunt motors, SCR crowbars, computer logic circuits, timing circuits, and inverters.
Comparison with SCS:
A silicon-controlled switch (SCS) behaves nearly the same way as an SCR; but there are a few differences. Unlike an SCR, an SCS switches off when a positive voltage/input current is applied to another anode gate lead. Unlike an SCR, an SCS can be triggered into conduction when a negative voltage/output current is applied to that same lead.
SCSs are useful in practically all circuits that need a switch that turns on/off through two distinct control pulses. This includes power-switching circuits, logic circuits, lamp drivers, and counters.
Compared to TRIACs:
A TRIAC resembles an SCR in that both act as electrically controlled switches. Unlike an SCR, a TRIAC can pass current in either direction. Thus, TRIACs are particularly useful for AC applications. TRIACs have three leads: a gate lead and two conducting leads, referred to as MT1 and MT2. If no current/voltage is applied to the gate lead, the TRIAC switches off. On the other hand, if the trigger voltage is applied to the gate lead, the TRIAC switches on.
Compared to TRIACs:
TRIACs are suitable for light-dimming circuits, phase-control circuits, AC power-switching circuits, AC motor control circuits, etc. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Endodontics**
Endodontics:
Endodontics (from the Greek roots endo- "inside" and odont- "tooth") is the dental specialty concerned with the study and treatment of the dental pulp.
Overview:
Endodontics encompasses the study (practice) of the basic and clinical sciences of normal dental pulp, the etiology, diagnosis, prevention, and treatment of diseases and injuries of the dental pulp along with associated periradicular conditions.
Overview:
In clinical terms, endodontics involves either preserving part, or all of the dental pulp in health, or removing all of the pulp in irreversible disease. This includes teeth with irreversibly inflamed and infected pulpal tissue. Not only does endodontics involve treatment when a dental pulp is present, but also includes preserving teeth which have failed to respond to non-surgical endodontic treatment, or for teeth that have developed new lesions, e.g., when root canal re-treatment is required, or periradicular surgery.Endodontic treatment is one of the most common procedures. If the dental pulp (containing nerves, arterioles, venules, lymphatic tissue, and fibrous tissue) becomes diseased or injured, endodontic treatment is required to save the tooth.
Procedures:
Root canal treatment Root canal treatment is a dental procedure used to treat infected tooth pulp which would be otherwise extracted. The pulp is the soft tissue core of the tooth which contains nerves, blood supply and connective tissue necessary for tooth health. This is usually caused when bacteria enter the pulp through a deep cavity or failed filling.Root canal treatment is required when the dental pulp is irreversibly damaged and involves both coronal and apical pulp. Root canal treatment can also be carried out on teeth with doubtful pulpal state before placing post-retained crowns and overdentures. Root canal therapy is not only performed when pain relief from an infected or inflamed pulp is required. It is also done to prevent adverse signs and symptoms from the surrounding sequelae and promote the healing and repair of the surrounding periradicular tissues. An example of which is if there is trauma to a front tooth which has caused it to be avulsed from the bony socket; endodontic treatment is required following re-implantation to preserve the aesthetics and function of the tooth, even though there may be no adverse symptoms of the dental pulp, or pain present at the time.Prior to root canal treatment, clinical examination and radiographic examinations are carried out to diagnose and plan treatment. Local anaesthesia is delivered to make the procedure pain free. The tooth to be treated is then isolated using a rubber dam, which prevents saliva entering the tooth during treatment and protects the airway from the fine files and strong chemicals used. The root canal treatment procedure is often carried out over single or multiple appointments. Root canal treatment involves: Removing the damaged and infected pulp Shaping the entire root canal system Cleaning and disinfecting the entire root canal system Filling and sealing the root canal system Placing a direct restoration such as composite filling or indirect restoration such as a crownInstrument fractures are common procedural mishaps in root canal treatment. It is essential to prevent instrument fractures. The success of removing of broken instrument depends on the location, direction and type of instrument. A sodium hypochlorite accident can result in long-term functional and aesthetic complications. Extrusion of sodium hypochlorite irrigating solution during a root canal procedure can cause a severe inflammatory reaction and tissue damage. Treatment is provided based on the severity of the injury. Tooth discolouration as a result of root canal treatment can occur if the pulpal tissue remnants are not completely removed or if a root canal sealer material containing silver is used.
Procedures:
Periradicular surgery Periradicular surgeries involve the root surface. These include apicoectomy (removal of a root end), root resection (removal of an entire root), repair of an injured root due to perforation or resorption, removal of broken fragments of the tooth or a filling material, and exploratory surgery to look for root fractures.
Apicoectomy An apicoectomy is a surgical procedure through which the apex of a root is resected, and a root-end filling is placed, preventing bacterial leakage into the root canal system from the periradicular tissues. A microsurgical technique is used to carry out apicectomy, which improves post-operative healing.
An apicoectomy can be carried out when a previous root canal treatment fails, and re-root canal treatment is not possible. This may be as a result of anatomical features, such as root dilaceration, which can compromise the completion of cleaning and obturating the root canal system. Procedural errors including ledges or perforations, may also be indications for an apicectomy.
Procedures:
Local anaesthetic is utilised to achieve anaesthesia as well as haemostasis for improved visualisation. A flap in the gum is designed, and then raised to allow for exposure of the periapical lesion. Bone removal (osteotomy) is carried out to enable access to root apex, and diseased tissue is removed at this point through curettage. The root end resection is carried out, removing 3mm apically. The canal(s) is then obturated, and the flap is sutured. There are a number of root-end filling materials available, including zinc oxide eugenol cements, and mineral trioxide aggregate.
Procedures:
Complications that may arise include: pain: anti-inflammatory agents or analgesics should be taken swelling: intermittent ice will aid in eradicating this. Swelling resolves usually within 24–48 hours.
Procedures:
ecchymosis (discolouration): this will often occur distant from the surgical site paraesthesia: usually transient as a result of inflammatory swelling, and sensation will return to normal in 4 weeks serious infection is rare, but can be treated with antibiotics, which should be administered with caution to avoid bacterial resistance maxillary sinus perforation Other procedures Other non-surgical endodontic procedures include pulp capping, pulpotomy, apexification, and pulpal regeneration. Hemisection, where a root and its overlying portion of the crown are separated from the rest of the tooth and optionally removed, is another (non-periradicular) endodontic surgery.
Tools:
Microsurgical endodontics, the use of magnification devices such as microscopes, and dental loupes, has been widely accepted among endodontists and practitioners; its use is believed to increase accuracy and visualization in the operating field. However, a Cochrane review in 2015 found no evidence to determine whether there is a difference in the outcome of a procedure done by magnification devices or a conventional procedure done with no magnification. The American Association of Endodontists strongly encourages its members to pursue the use of an oral microscope to ensure the highest level of excellence.The use of a CBCT is also becoming the standard of care.
Training:
Endodontists are specialist dentists with additional training, experience and formal qualifications in endodontic treatment, apicectomies, microsurgery, and dental emergency and trauma management. Endodontics is recognized as a specialty by many national dental organizations including the Dental Board of Australia, British General Dental Council, American Dental Association, Royal College of Dentists of Canada, Indian Dental Association, and Royal Australasian College of Dental Surgeons.
Training:
Australia In Australia, endodontics is recognized as one of the thirteen registered dental specialties. In addition to a dental degree, Endodontists have an additional three years of postgraduate University training in the area of Endodontics to be recognized and registered by the Dental Board of Australia. A general dentist is permitted by law to perform endodontic treatment, but must be competent in the skills required for the endodontic treatment, and refer complex cases for specialist management.
Training:
United States In the United States after finishing a dental degree, a dentist must undergo 2–3 additional years of postgraduate training to become an Endodontist. American Dental Association (CODA) accredited programs are a minimum of two years in length. Following successful completion of this training, the dentist becomes Board eligible to sit for the American Board Of Endodontology examination. Successful completion of board certification results in Diplomate status in the American Board of Endodontics.Although general dentists can perform endodontic treatment, there are several things which set endodontists apart. Endodontists use microscopes during treatment to better treat the small internal anatomy of teeth without taking away too much tooth structure, or causing iatrogenic damage. Also, endodontists use CBCT (3D imaging) to assess case difficulty and for diagnosis and treatment planning of endodontic cases. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Solar energy in Finland**
Solar energy in Finland:
Solar energy in Finland is used primarily for water heating and by the use of photovoltaics to generate electricity. As a northern country, summer days are long and winter days are short. Above the Arctic Circle, the sun does not rise some days in winter, and does not set some days in the summer. Due to the low sun angle, it is more common to place solar panels on the south side of buildings instead of on the roof. Mounting them vertically reduces the average output by 22% from mounting at a 60° angle.
Photovoltaics:
The PV capacity of Finland was (2012) 11.1 MWp. Solar power in Finland was (1993–1999) 1 GWh, (2000–2004) 2 GWh and (2005) 3 GWh. There has been at least one demonstration project by the YIT Rakennus, NAPS Systems, Lumon and City of Helsinki in 2003. Finland is a member in the IEA's Photovoltaic Power Systems Programme but not in the Scandinavian Photovoltaic Industry Association, SPIA.
Photovoltaics:
In 2015, the Kaleva Media printing plant in Oulu became the most powerful photovoltaic solar plant in Finland, with 1,604 solar photovoltaic (PV) units on its roof. Although the city of Oulu, located near the Arctic Circle, has only two hours of weak sunlight in December, the photovoltaic cells work almost around the clock in the summer. The cold climate means the PV panels can get up to a 25% boost per hour, as they don't overheat.Because the sun is quite low in the sky at this latitude, vertical PV installations are popular on the sides of buildings. These solar walls also capture light reflected from snow.Snow is not necessarily cleared from rooftop solar installations.
Solar heating:
The objective in solar heating is 163 000 m2 collector area (1995–2010). In 2006 the collector area in operation was 16 493 m2. Solar heat in Finland was (1997–2004) 4-5 GWh and (2005) 6 GWh. Thus, Finland has installed 10% of its objective in 11 years time (1995–2010). The solar heating has not been competitive due to cheap alternatives (electricity, fuel oil and district heating) and the lack of support systems. Companies and public organizations may receive 40% investment subsidies, but private houses do not receive subsidies yet. The Finnish Solar Industries (FSI) group was established in 2001. 2006/2005 the markets grew 43%. Finland's production capacity is 16 000 m2/a. New installations were: 2 380 m2 (2006), 1 668 m2 (2005) and 1 141 m2 (2004). There are growth opportunities in the solar heating. In 2018 S-Ryhmä decided to order solar panels for 40 of its commercial real estate buildings. This is the biggest solar panel project in Finnish history. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Krauss-Helmholtz bogie**
Krauss-Helmholtz bogie:
A Krauss-Helmholtz bogie (Krauss-Helmholtz-Lenkgestell) is a mechanism used on steam locomotives and some electric locomotives to improve curve running.
Operation:
The bogie comprises a carrying axle connected to a coupled axle via a shaft or lever. In straight running, any radial movement of the carrying axle results in a sideways movement of the coupled axle in the opposite direction. However, the carrying axle is centred by means of two heavy duty springs just behind it. In addition the pivot pin may be allowed to move sideways, but again is held centrally by heavy springs. When travelling round a curve, the carrying axle swings to one side causing the coupled axle to move sideways in the opposite direction. In this way radial forces during curve running are more or less evened out on both axles, so that riding qualities similar to those of a normal bogie are achieved and wear and tear reduced on wheel flanges and rails.
Operation:
The bogie is a type of pony truck and was named after the locomotive firm of Krauss and the engineer, Richard von Helmholtz. By contrast a Bissel bogie is independently installed in the frame, and sideways guidance of the locomotive is achieved by elastic forces. The distribution of these forces is not tightly defined and, in addition, they are dependent on the curve radius.
Examples of use:
Steam locomotives Because the advantages of a pony truck come into play particularly on tight curves, the Krauss-Helmholtz bogie initially appeared on branch line, Lokalbahn and narrow gauge locomotives. One of the first locomotives of this type was the Bavarian D VIII. On this tank locomotive the bogie was located at the rear; however in the majority of cases it was at the front or - if the locomotive had to have equally good riding qualities in both directions - at both ends.
Examples of use:
Later this pony truck arrangement was also adopted by the DRG standard locomotives (Einheitslokomotive) of the Deutsche Reichsbahn, e.g. on the ten-coupled classes: 44, 45, 50 and 85. An exception was the Class 84, that was fitted with Schwartzkopff-Eckhardt II bogies or Luttermöller axle drives.
The tender locomotives of classes 41 and 45 only had a Krauss-Helmholtz bogie at the front; the trailing axle was housed in a Bissel bogie. The tank locomotives of Class 85, like some of the Class 64 and 86 engines, had two Krauss-Helmholtz bogies.
Electric locomotives Even the electric locomotives of Reichsbahn classes E 04, E 17, E 18 and E 19 were fitted with comparable pony trucks, known as AEG frames (AEG-Gestell). Because the axles had external bearings, the lever linkage also had to be on the outside, a characteristic detail of these locomotives.
Examples of use:
Italian bogie In Italy, the Krauss-Helmholtz bogie was improved around 1900 by Giuseppe Zara, a technician of the Rete Adriatica and later of the Ferrovie dello Stato; by modifying the structure, rearranging the weight distribution and allowing the bogie to also move transversally respective to the locomotive's frames, he obtained the Italian bogie (a lighter and somewhat different version was the later Zara bogie). It saw widespread use, being fitted on many Italian steam locomotive classes, like the FS Class 625, FS Class 740 and FS Class 685. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Multi-spectral phase coherence**
Multi-spectral phase coherence:
Multi-spectral phase coherence (MSPC) is a generalized cross-frequency coupling metric introduced by Yang and colleagues in 2016. MSPC can be used to quantify nonlinear phase coupling between a set of base frequencies and their harmonic/intermodulation frequencies. MSPC is a model-free method, which can provide a system description, including (i) the order of the nonlinearity, (ii) the direction of interaction, (iii) the time delay in the system, and both (iv) harmonic and (v) intermodulation coupling.
Multi-spectral phase coherence:
The MSPC is defined as: exp sum )))⟩ where φ(fi) is the phase at frequency fi , ai is the weight of fi to a harmonic/intermodulation frequency sum =∑iaifi ), and ⟨⋅⟩ represents the average over realizations.
Bi-phase locking value, also called bi-phase coherence in the literature, is a special case of MSPC when a1=a2=1 , 2.
The time-delay can be estimated from the phase lag when MSPC is computed between signals. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**TRNA (cytosine-5-)-methyltransferase**
TRNA (cytosine-5-)-methyltransferase:
In enzymology, a tRNA (cytosine-5-)-methyltransferase (EC 2.1.1.29) is an enzyme that catalyzes the chemical reaction S-adenosyl-L-methionine + tRNA ⇌ S-adenosyl-L-homocysteine + tRNA containing 5-methylcytosineThus, the two substrates of this enzyme are S-adenosyl methionine and tRNA, whereas its two products are S-adenosylhomocysteine and tRNA containing 5-methylcytosine.
This enzyme belongs to the family of transferases, specifically those transferring one-carbon group methyltransferases. The systematic name of this enzyme class is S-adenosyl-L-methionine:tRNA (cytosine-5-)-methyltransferase. Other names in common use include transfer ribonucleate cytosine 5-methyltransferase, and transfer RNA cytosine 5-methyltransferase. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Serrated blade**
Serrated blade:
A serrated blade is one with a toothlike rather than a plain edge, and is used on saws and on some knives and scissors. It is also known as a dentated, sawtooth, or toothed blade. Most such blades are scalloped, having edges cut with curved notches, common on wood saws and bread knives. Serrations give the blade's cutting edge less contact area than a smooth blade, which increases the applied pressure at each point of contact, and the points of contact are at a sharper angle to the material being cut. This causes a cutting action that involves many small splits in the surface of the material being cut, which cumulatively serve to cut the material along the line of the blade.Cuts made with a serrated blade are typically less smooth and precise than cuts made with a smooth blade. Serrated edges can be difficult to sharpen using a whetstone or rotary sharpener intended for straight edges but can be sharpened with ceramic or diamond coated rods. Further, they tend to stay sharper longer than similar straight edges. A serrated blade has a faster cut, but a plain edge has a cleaner cut. Some prefer a serrated blade on a pocket knife.
Types of serration:
Tooth serration — Vertical serration along edge of blade Single edge serration — Serration on one side, the other remains flat Double edge serration — Serration on both sides Fan serration — Side-to-side serration without necessarily having a toothed edge Micro-serration — Serration much smaller than thickness of blade creating something like a fan pattern | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Isopropylphenidate**
Isopropylphenidate:
Isopropylphenidate (also known as IPH and IPPD) is a piperidine based stimulant drug, closely related to methylphenidate, but with the methyl ester replaced by an isopropyl ester. It has similar effects to methylphenidate but with a longer duration of action, and was banned in the UK as a Temporary Class Drug from April 2015 following its unapproved sale as a designer drug.It has been researched as potential methylphenidate replacement for ADHD and narcolepsy, because of fewer side effects. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Indium trihydride**
Indium trihydride:
Indium trihydride is an inorganic compound with the chemical formula (InH3). It has been observed in matrix isolation and laser ablation experiments. Gas phase stability has been predicted. The infrared spectrum was obtained in the gas phase by laser ablation of indium in presence of hydrogen gas InH3 is of no practical importance.
Chemical properties:
Solid InH3 is a three-dimensional network polymeric structure, where In atoms are connected by In-H-In bridging bonds, is suggested to account for the growth of broad infrared bands when samples of InH3 and InD3 produced on a solid hydrogen matrix are warmed. Such a structure is known for solid AlH3. When heated above −90 °C, indium trihydride decomposes to produce indium–hydrogen alloy and elemental hydrogen. As of 2013, the only known method of synthesising indium trihydride is the autopolymerisation of indane below −90 °C.
Other indium hydrides:
Several compounds with In-H bonds have been reported. Examples of complexes with two hydride ligands replaced by other ligands are (K+)3[K((CH3)2SiO)+7][InH(CH2C(CH3)3)−3]4 and HIn(−C6H4−ortho-CH2N(CH3)2)2.
Although InH3 is labile, adducts are known with the stoichiometry InH3Ln (n = 1 or 2).
1:1 amine adducts are made by the reaction of Li+[InH4]− (lithium tetrahydridoindate(III)) with a trialkylammonium salt. The trimethylamine complex is only stable below −30 °C or in dilute solution. The 1:1 and 1:2 complexes with tricyclohexylphosphine (PCy3) have been characterised crystallographically. The average In-H bond length is 168 pm. Indium hydride is also known to form adducts with NHCs. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Anti-nesting principle**
Anti-nesting principle:
In the philosophy of consciousness, the anti-nesting principle states that one state of consciousness cannot exist within another.Proponents of the anti-nesting principle include Giulio Tononi and Hilary Putnam. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Simultaneously extracted metals and acid-volatile sulfide**
Simultaneously extracted metals and acid-volatile sulfide:
Simultaneously extracted metals/Acid-volatile sulfide (SEM-AVS) is an approach used in the field of aquatic toxicology to assess the potential for metal ions found in sediment to cause toxic effects in organisms dwelling in the sediment. In this approach, the amounts of several heavy metals in a sediment sample are measured in a laboratory; at the same time, the amount of acid-volatile sulfide (sulfide which can be liberated from the sediment by treatment with hydrochloric acid) is determined. Based on the chemical interactions between heavy metals (SEM) and acid-volatile sulfide (AVS), the concentrations of these two components can be used to assess the potential for toxicity to sediment-dwelling organisms.
Background:
Metals A number of heavy metals, such as cadmium and lead, are toxic to various forms of life, particularly when dissolved in water as metal ions. Toxic heavy metals are often present in surface water as a result of natural processes, such as the weathering of metal-containing rocks, or due to human activity, such as mining and smelting. Only the ionic forms of most metals are soluble in water. These ionic forms have a high chemical affinity for the surfaces of most sediment particles, meaning they bind tightly to their surface. As a result, when water bearing heavy metal ions is in contact with sediment, the ions tend to accumulate in the sediment at elevated concentrations. This is an example of partition equilibrium. If metal ions are present in great enough quantities, they may have toxic effects on organisms that are exposed to them by ingestion or absorption.
Background:
Sulfide The sulfide ion (S2−) is present in some anoxic sediments as a result of bacterial activity. In environments containing little or no oxygen gas (O2) but large amounts of sulfate ion (SO42−), sulfate-reducing bacteria use sulfate in their metabolism as an electron acceptor. This process creates sulfide as a product according to Equation 1.
The sulfide ion produced by this process is sensitive to biological or chemical oxidation in the presence of oxygen, so it only persists in sediments that are continuously anoxic.
Background:
Metal-Sulfide Interactions When dissolved in water, sulfide has a high affinity for numerous heavy metal ions. That is, the solubility-product constants (Ksp) for the sulfides of these metals are very small, meaning they will precipitate as solids when a heavy metal ion and sulfide ion come into contact, as in Equation 2, where M2+ is a generic metal in the +2 oxidation state.
Background:
In anoxic sediments uncontaminated by heavy metals, the associated metal (M in equation 2) is usually iron (Fe) or manganese (Mn). Iron (II) is abundant in anoxic sediment, and the Ksp for iron (II) sulfide is 10−22.39 (with a comparable value for manganese (II) sulfide), so effectively all the sulfide in an uncontaminated sediment will be bound to Fe or Mn.
Background:
Several toxic heavy metals, however, have Ksp values substantially lower than those of the sulfides of iron and manganese - for example, cadmium (II) sulfide (CdS) has a Ksp equal to 10−32.85. This means cadmium binds sulfide with a much higher affinity than does iron. When water contaminated with cadmium ions comes into contact with sediment containing FeS, the cadmium displaces the iron according to Equation 3 and becomes tightly bound to the sulfide ion.
Background:
Due to the large difference in Ksp values for the two metal sulfides, this reaction proceeds effectively to completion, meaning that until all the sulfide in a sediment is used up, all the cadmium in that sediment will be present in the solid CdS form.
A number of other toxic heavy metals behave similarly, including lead, copper, zinc, mercury, and nickel.
Background:
Bioavailability In order for toxic substances like heavy metals to cause effects in organisms, they must be bioavailable. For organisms residing in contaminated sediments, the contaminants are most bioavailable when dissolved in the pore water, as opposed to being precipitated as a solid or sorbed to a sediment particle. Metals in the solid metal-sulfide form are thus considered non-bioavailable, and are unlikely to cause toxicity in sediment-dwelling organisms.
Background:
Thus, sediments with the same quantity of metals in them may have vastly different toxic effects, depending on the quantity of sulfide available to bind with them and render them non-bioavailable. For this reason, the SEM-AVS approach was developed to account for differences in sulfide and refine methods for predicting heavy metal toxicity in sediments.
Methods:
Sample Collection Because sulfide is quickly degraded in the presence of oxygen, sediment samples for SEM/AVS analysis must be kept under rigorously anoxic conditions from the moment they are sampled. In addition, samples should be kept at 4 °C to keep bacterial metabolism from altering sediment composition. The State of Ohio Environmental Protection Agency recommends storing samples for no longer than 14 days before analyzing them.
Methods:
Extraction Sediment samples to be analyzed are first purged with argon or nitrogen gas to ensure they are anoxic. The sample is placed in a flask connected to an apparatus for trapping hydrogen sulfide gas (H2S). Oxygen-free water and hydrochloric acid (HCl) are added, and the sediment is stirred for one hour while argon or nitrogen gas is bubbled through.
Methods:
Sulfide Determination When HCl is mixed with metal sulfides in the sediment, a reaction occurs that generates H2S and liberates the metal ion into aqueous form, as shown in Equation 4. The gas formed by this process accumulates in the trap connected to the flask. By weighing the trap before and after the extraction process, the amount of H2S produced by the reaction can be calculated.
Methods:
A few important things should be noted about this reaction: Not all of the sulfide present in a sediment sample will undergo this reaction when exposed to HCl. Some may be buried deep inside sediment grains and be unavailable to react. Thus, the method is a measurement of the "acid-volatile sulfide," rather than the total sulfide.
The stoichiometric ratio of MS to H2S is 1:1, meaning for every sulfide ion present as a metal sulfide compound, one molecule of H2S is generated. Thus, the amount (in moles) of H2S measured by weighing the trap is equal to the amount of AVS originally present in the sediment.
Methods:
The metal (M) is transformed from solid form to the dissolved form. In uncontaminated sediment, M will mostly be a combination of Fe and Mn. In contaminated sediments, toxic heavy metals such as Cd, Pb, etc. will also be liberated by the reaction.Once the quantity (in moles) of AVS has been determined in this way, it is divided by the dry mass of the sediment to obtain the AVS concentration. In addition to the gravimetric method described here, other methods, such as colorimetry, may be used.
Methods:
Metals Determination As noted above, treating metals-containing sediments with HCl liberates metal ions into the acid solution that were previously bound up with AVS. After treatment these are present in solution, along with any metals that were initially unbound to AVS (due to insufficient AVS in the sediment to "mop up" all the metal ions). The concentrations of metals in the acid solution can be measured by a number of analytical chemistry methods, including atomic emission spectrometry, atomic absorption spectrometry (AAS), or mass spectrometry. These are known as "simultaneously extracted metals" because they are the metals liberated from the sediment while the volatilization of sulfide is occurring. Metals extracted from sediment by digestion, or with a different acid than HCl are not simultaneously extracted metals. By correcting appropriately for dilution, the SEM concentration in the sediment can be determined.
Methods:
In initial versions of the SEM-AVS approach, six metals were measured: nickel, zinc, cadmium, copper, lead, and mercury. More current methods call for the measurement of silver in addition to these metals.
Toxicity:
Theory In applying the SEM-AVS approach, two concentrations are determined: the total concentrations of all toxic heavy metals of interest, represented as [SEM], and the acid-volatile sulfide concentration, represented as [AVS]. From these concentrations, the [SEM]/[AVS] ratio can be obtained, summarizing the results in a single value.
Toxicity:
Based on the extremely low Ksp values for heavy metal sulfides, if [SEM]/[AVS] is less than 1 ([SEM < [AVS]), then all the extractable metals in the sediment should be in their solid sulfide form and non-bioavailable; above this value, the pool of sulfide is "exhausted" and heavy metals are more likely to be present in the sediment pore water, their bioavailable form. In theory, then, an [SEM]/[AVS] value of 1 represents a cutoff between a "no-effects" range and a "possible effects" range.
Toxicity:
Spiked-Sediment Studies A 1996 study of the predictive power of the [SEM]/[AVS] approach employed laboratory toxicity testing of spiked sediments. Toxic heavy metals, alone and in combination, were added to clean sediments with varying concentrations of AVS. Benthic organisms were then exposed to the sediments, and their mortality was measured and compared to metal-free controls. Ninety-two different trials were conducted, using several test species exposed to cadmium, copper, nickel, lead, and zinc.
Toxicity:
In sediments where [SEM]/[AVS] was less than or equal to 1, only 1.1% of trials showed greater toxicity than in controls. Where the ratio was greater than 1, 73.5% of trials showed greater toxicity than controls. These results held for both fresh and salt water, for different metal types and test species, and across a range of SEM and AVS concentrations. A study that deployed spiked sediments in a pond found a similar threshold at [SEM]/[AVS] = 1 for effects to local benthic fauna. Spike studies that measured heavy metals concentrations in the pore water - that is, before extracting with HCl - found that, when [AVS] was greater than [SEM], pore water heavy metal concentrations were undetectable or nearly so. These results strongly support the basic theoretical framework of the SEM-AVS approach.
Toxicity:
Field Studies The results of applying the SEM-AVS approach to contaminated sites in the field were mixed. One study found that [SEM]/[AVS] was a good predictor of toxic effects in laboratory exposures to field-contaminated sites. However, a re-analysis by other investigators suggested that [SEM]/[AVS] was neither a more sensitive predictor of toxic effects nor a more efficient one than simply measuring metals concentration per sediment dry weight. In addition, in Flemish rivers polluted with metals, AVS concentration had little to no effect on the accumulation of metals in benthic organisms (though no measurements of toxicity were done).
Sediment Quality Assessment:
The potential for SEM-AVS to act as a screening tool in evaluating sediment toxicity due to metals has led several regulatory bodies to use it in the establishment of regulatory assessments of sediment quality.
European Union Under the Water Framework Directive, implemented by the European Union in 2003, Environmental Quality Standards (EQS) for metals-contaminated sediments incorporate [SEM]/[AVS] measurements.
Australia In sediment quality assessments in Australia, SEM-AVS is employed in second-tier assessment. That is, when initial screening indicates metals concentrations in excess of guidelines, [SEM]/[AVS] is calculated to determine if sufficient sulfide is present to mitigate metal bioavailability.
Sediment Quality Assessment:
United States In a 2002 guidance manual for assessment of contaminated freshwater sediment, the US Environmental Protection Agency (EPA) listed [SEM] - [AVS] (a variant of the approach in which data are normalized by subtraction rather than division) as a "Moderate Priority" metric of ecosystem health. Since 2005, EPA Equilibrium Partitioning Sediment Benchmarks (ESBs) for metal mixtures have been derived from a combination of Water Quality Criteria (WQC) and a variant of [SEM]/[AVS] corrected for the organic carbon content of the sediment.
Sediment Quality Assessment:
In conducting its National Sediment Quality Survey, the EPA used [SEM] - [AVS] in its classification scheme; a sediment with an [SEM] - [AVS] concentration greater than 5 μg/g dry weight was classified as "Tier 1: Associated Adverse Effects on Aquatic Life or Human Health Are Probable." Values between 0 and 5 μg/g dry weight were classified as "Tier 2: Associated Adverse Effects on Aquatic Life or Human Health Are Possible." The Ohio Environmental Protection Agency uses [SEM]/[AVS] in sediment quality guidelines for waters in that state. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Medium of instruction**
Medium of instruction:
A medium of instruction (plural: media of instruction, or mediums of instruction) is a language used in teaching. It may or may not be the official language of the country or territory. If the first language of students is different from the official language, it may be used as the medium of instruction for part or all of schooling. Bilingual education or multilingual education may involve the use of more than one language of instruction. UNESCO considers that "providing education in a child's mother tongue is indeed a critical issue". In post-secondary, university and special education settings, content may often be taught in a language that is not spoken in the students' homes. This is referred to as content based learning or content and language integrated learning (CLIL). In situations where the medium of instruction of academic disciplines is English when it is not the students' first language, the phenomenon is referred to as English-medium instruction (EMI).
In different countries and regions:
Africa In Tanzania, Swahili is used in primary schools and adult education, whereas English is used in secondary schools and universities.
In Zimbabwe, use of English, Shona and Ndebele is established in education until the fourth grade; from the fourth grade, English is the medium of instruction.
In different countries and regions:
In South Africa, students are taught primarily in their home language from Grade Zero (Reception Year) up to Grade 3. From Grade 4 onwards, English is the default language of learning and teaching, except for a minority of schools in which Afrikaans is used. The national curriculum requires that all students study at least two official languages as separate subjects, one of which must be studied at home language level and the other at least at first additional language level. The most common home language among the school population is isiZulu.
In different countries and regions:
In Nigeria, Medium of Instruction at all levels of education (primary, secondary, universities and colleges) is English Language.
In the francophone states of Africa, education has typically been in French only.
In Ethiopia, Amharic, Oromo, and other Ethiopian languages serve as the medium of instruction in primary education, while English is used in secondary schools and universities (French had been the medium of instruction in public schools pre-1936).
In different countries and regions:
Western Hemisphere Brazil Every public school uses Brazilian Portuguese as the medium of instruction, but no law prohibits the use of other languages in private schools. Many schools use other European languages (mainly because of the country's European heritage) such as English, German, Italian or French. Public schools also have mandatory English and Spanish but only once or twice a week.
In different countries and regions:
Canada In Canada, almost all public schools use either English or French as the medium of instruction; French is standard in the province of Quebec (a few cities also offering English-language schools) and, along with English, in New Brunswick. The official language not used as the primary medium of instruction is taught as a mandatory subject in primary school and becomes optional for most secondary school students. Many public and private school systems in English jurisdictions also offer French immersion. In parts of Canada such as the Northwest Territories and Nunavut, some aboriginal languages, such as Inuinnaqtun and Tłı̨chǫ, are also used in local school systems. Heritage languages (immigrant or diaspora languages) are as also common across Canada in public and private schools.
In different countries and regions:
United States English is used, but in some schools, Spanish, French (in Louisiana), Hawaiian (in Hawaii), and local Native American/American Indian languages are used as well.
In different countries and regions:
The Cherokee Nation instigated a 10-year language preservation plan that involved growing new fluent speakers of the Cherokee language from childhood on up through school immersion programs as well as a collaborative community effort to continue to use the language at home. This plan was part of an ambitious goal that in 50 years, 80% or more of the Cherokee people will be fluent in the language. The Cherokee Preservation Foundation has invested $3 million into opening schools, training teachers, and developing curricula for language education, as well as initiating community gatherings where the language can be actively used. Formed in 2006, the Kituwah Preservation & Education Program (KPEP) on the Qualla Boundary focuses on language immersion programs for children from birth to fifth grade, developing cultural resources for the general public and community language programs to foster the Cherokee language among adults. There is also a Cherokee language immersion school in Tahlequah, Oklahoma that educates students from pre-school through eighth grade.
In different countries and regions:
Asia In Azerbaijan, Azeri is the main language of instruction. Instruction in Russian and to a lesser extent in English at both secondary and postsecondary level is also offered in some educational institutions. Georgian is the language of instruction in secondary schools in the Georgian-populated northern regions. Education was conducted in Armenian in some secondary schools until the First Nagorno-Karabakh War.
In different countries and regions:
In Bangladesh, Bengali and English are used as mediums of instruction. In universities, the medium of education is mainly English.
In the People's Republic of China (PRC), Standard Chinese is used as the medium of instruction in most schools. In elementary and secondary schools for ethnic minorities, the minority languages, such as Mongolian, Uyghur, Tibetan and Korean, are also used.
In Georgia, most schools conduct education in Georgian. The number of Azerbaijani schools is being reduced.
In different countries and regions:
In Hong Kong which was a British colony until 1997, either English or Cantonese is the medium in most schools at the primary and secondary level. EMI schools adopt English as the medium of instruction for almost all classes. CMI schools generally adopt Cantonese as the medium of instruction, but a significant number of CMI schools use English for high school courses. English is used almost exclusively at the tertiary level.
In different countries and regions:
In Israel, Hebrew is the medium in most schools, and Arabic is the medium in elementary and secondary schools for the Arab minority. Hebrew is used almost exclusively at the tertiary level.
In different countries and regions:
In India, the medium of instruction varies among English, Hindi and the respective states’ official languages. Private schools usually prefer English, and government (primary/secondary education) schools tend to go with one of the last two. However, the medium of instruction in colleges and universities is always either English, Hindi or a regional language. The medium of education is also dependent upon the state and its official language.
In different countries and regions:
In Assam, Assamese or English is used.
In West Bengal, Bengali or English is used.
In Karnataka, Kannada or English is used.
In Goa, English or Konkani is used.
In Gujarat, Gujarati or English is used.
In Maharashtra, English or Marathi is used.
In Andhra Pradesh, Telugu or English is used. Some schools also use Sanskrit.
In Telangana, Telugu, or English is used.
In Tamil Nadu, Tamil or English is used.
In Kerala, Malayalam or English is used.
In Odisha, Odia or English is used.
In Japan, Japanese is used in most schools (including universities and colleges).
In South and North Korea, Korean is used in most schools (including universities and colleges).
In Macau, Cantonese is used as the medium of instruction in many schools. Portuguese is used in Portugal-backed schools. English, which is not an official language of the region, is also used in a lot of schools.
In Pakistan, most public schools use Urdu, but private schools have English as medium of instruction. English was made medium of instruction in 18 colleges in 2008.Government of the Punjab Notification No. PS/SSE/Misc/2009/176 dated 18-09-2009 required for the subjects of science and mathematics to be taught in English in each school.
In Taiwan, Standard Chinese is used as the medium of instruction in most schools, with more Taiwanese Hokkien presence in schools in recent years.
South East Asia In the Philippines, English is the primary medium of instruction from preschool to university, except in the Philippine history and Filipino language subjects, in which Filipino is used. Recently, regional languages have been introduced as the medium of instruction in public schools for grades K–3 as part of the Department of Education's mother tongue-based education policy.
In Singapore, in pre-schools children learn in two languages: English and a mother tongue: Chinese, Malay or Tamil. The medium of instruction is English in all schools following the national curriculum except in "mother-tongue" subjects. International and private schools may use other languages. See also Special Assistance Plan.
In Indonesia, Indonesian is the medium in most schools, including universities.
In Cambodia, Khmer is the medium in most schools, including universities.
In Vietnam, Vietnamese is the medium in most schools, including universities.
In Thailand, Thai is the medium in most schools, including universities.
In Laos, Lao is the medium in most schools, including universities.
In different countries and regions:
In Malaysia, Malay is the medium of instruction in most schools. However, there are also Chinese and Indian schools serving the respective communities, which are allowed to use Mandarin and Tamil respectively as a medium of instruction, but Malay is still required to be taught as a subject. English-medium schools were present during the colonial period but were slowly phased out after independence. Today, all the former English-medium schools have since been converted to Malay-medium schools. Nevertheless, English continues to be a compulsory subject in all Malaysian schools.
In different countries and regions:
Australia and Oceania In Australia, most schools use English. However, in the State of Victoria (known for its many Greek and Italian settlers), there are a number of schools that teach in Greek and Italian. A few schools also teach in French, Chinese, Arabic and Japanese.
In New Zealand, English is used in many schools, but more and more kohanga reo (kindergarten) and kura kaupapa (primary and secondary school) are using Māori instead.
In Vanuatu, English and French are the main languages of education.
Europe In Belarus, Russian is the main language of instruction. While schools using Belarusian schools are 53%, they are located mostly in rural areas, and the share of students who receive instruction in Belarusian is as low as 18%.
In Belgium, Dutch and French (German in some parts of Eastern Belgium) are used.
In Croatia, besides Croatian-language education, education of the representatives of national minorities is carried out in 24 elementary schools, and the program is conducted in the language and writing of a relevant national minority, 61 elementary schools having classes with such programs.
In Estonia, as of 2011, there were 463 Estonian-medium schools, 62 Russian-medium schools and 36 mixed medium schools, 25% of vocational education being in Russian while the remainder Estonian. In higher education, 90.2% is in Estonian, 7.8% in Russian and 1.85% in English.
In different countries and regions:
In Finland, Finnish is the language used in most schools, but Swedish, also an official national language, is used in a number of schools along the coast and Abo Akademi. The right to education in Swedish is based in the constitution. There are also a few schools in which education is given, to some extent, in Sami in the north. See also Mandatory Swedish.
In different countries and regions:
In France, legislation restricts languages other than French in state schools. Other languages of France are the medium of instruction in non-state schools such as Diwan Breton language-medium schools and the Calandretas in the south that use Occitan. See Language policy in France In Iceland, Icelandic is used at all levels of education. English is the first secondary language to be taught (even starting a bit as early as kindergarten), with Danish also required later. Some universities teach in part in English in topics popular with foreigners (and "Icelandic for foreign students" is also offered).
In different countries and regions:
In Ireland, English is used in most schools with a growing number of gaelscoileanna (10%) using Irish.
In Italy, while Italian is official language throughout the territory, also French is official in Valle D'Aosta, German in South Tirol.
In different countries and regions:
In Latvia, Latvian is used in most schools. According to the Ministry of Foreign Affairs, education is available in eight national minority languages: Russian, Polish, Ukrainian, Belarusian, Lithuanian, Estonian, Hebrew and Romani. Boris Tsilevitch, politician and former chair of PACE sub-commission on minorities, notes that all minority schools (except the Russian and Polish ones) offer education in either Latvian or Russian, with corresponding minority language and culture taught as subjects. The network of Russian-language schools is being reduced. Some Polish-language schools were created after restoration of independence. Education in public minority high schools is conducted mostly in Latvian since 2004 despite wide protests by the Headquarters for the Protection of Russian Schools.
In different countries and regions:
In Lithuania, as at 2004/2005, 91.3% of pupils studied in Lithuanian, 5.3% in Russian and 3.6% in Polish in general education schools.
In Moldova, Moldovan (Romanian) is used, but Russian is slowly being introduced.
In North Macedonia, the state is obliged by Ohrid Agreement to provide university level education in languages spoken by at least 20% of the population (Albanian) In Norway, the medium of instruction is Norwegian. The state undertakes to provide a substantial part of preschool education in Sami, at least pupils whose families request it in sufficient numbers.
In Poland, the medium of instruction in most schools is Polish. However, in areas where national or ethnic minorities reside, there are as well public schools with minority language of instruction (such as German, Ukrainian, Belarusian, or Kashubian), or schools which offers classes of the minority language.
In different countries and regions:
In Romania, the medium of instruction is mostly Romanian, but the state undertakes to provide education in minority languages up to the following levels. In Russian, it is a substantial part of preschool education, at least to those pupils whose families request it in sufficient numbers. In Bulgarian and Czech, it is a substantial part of primary education. In Croatian, it is a substantial part of secondary education. In Serbian, Turkish, Ukrainian and Slovak, it is secondary education. In German and Hungarian, it is higher education. There are also international schools where the medium of instruction is English.
In different countries and regions:
In Russia, Russian is dominating in education. Approximately 6% of students learn at school in minority languages. Besides, some tertiary education establishments use Tatar as a language of instruction alongside Russian.
In different countries and regions:
In Slovakia, education in minority languages must be provided in municipalities if Slovak citizens speaking respective language are more than 20% of population: higher, technical and vocational education in Hungarian, a substantial part of technical and vocational education in Ruthenian and Ukrainian, a substantial part of preschool education for those pupils whose families request it in sufficient numbers in Bulgarian, Croatian, Czech, German, Polish and Roma.
In different countries and regions:
In Slovenia, the general medium of instruction is Slovene. In areas with the Hungarian ethnic minority, bilingual instruction in Slovene and Hungarian is compulsory. In the Italian ethnic community area, basic education can be provided in Slovene or Italian.
In Switzerland, German, French, Italian and/or Romansh are used in most schools.
In different countries and regions:
In Ukraine since the 2017 law "On Education" is language of instruction in Ukrainian schools is the state language, which is Ukrainian (national minorities are guaranteed the right to study in public educational facilities including their language alongside Ukrainian).Prior to the 2017 law "On Education" the mediums of instruction in pre-school education were Ukrainian, Russian, Hungarian, Romanian, Moldovan, Crimean Tatar, English, Polish and German; in general education, Ukrainian, Russian, Hungarian, Romanian, Moldovan, Crimean Tatar, Polish, Bulgarian and Slovak; in vocational training, Ukrainian and Russian; in higher education, Ukrainian, Russian, Hungarian, Romanian.
In different countries and regions:
In the United Kingdom, English is mostly used.
In the Isle of Man, English is used, but Manx is being revived with one Manx-medium school at St. John’s.
In Scotland, English is the primary language of instruction but Gaelic medium education is also available. There is little or no use of Lowland Scots as a medium of education.
In Wales, most schools teach in English, but more and more teach in Welsh.
In Cornwall, Cornish is used in some schools that otherwise teach in English. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
Subsets and Splits
No community queries yet
The top public SQL queries from the community will appear here once available.