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python retrieve.py --class_prompt {} --class_data_dir {} --num_class_images 200 Then you can launch the script: Copied export MODEL_NAME="CompVis/stable-diffusion-v1-4" |
export OUTPUT_DIR="path-to-save-model" |
accelerate launch train_custom_diffusion.py \ |
--pretrained_model_name_or_path=$MODEL_NAME \ |
--output_dir=$OUTPUT_DIR \ |
--concepts_list=./concept_list.json \ |
--with_prior_preservation \ |
--real_prior \ |
--prior_loss_weight=1.0 \ |
--resolution=512 \ |
--train_batch_size=2 \ |
--learning_rate=1e-5 \ |
--lr_warmup_steps=0 \ |
--max_train_steps=500 \ |
--num_class_images=200 \ |
--scale_lr \ |
--hflip \ |
--modifier_token "<new1>+<new2>" \ |
--push_to_hub |
</hfoption> |
</hfoptions> |
Once training is finished, you can use your new Custom Diffusion model for inference. |
<hfoptions id="training-inference"> |
<hfoption id="single concept"> |
Copied import torch |
from diffusers import DiffusionPipeline |
pipeline = DiffusionPipeline.from_pretrained( |
"CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16, |
).to("cuda") |
pipeline.unet.load_attn_procs("path-to-save-model", weight_name="pytorch_custom_diffusion_weights.bin") |
pipeline.load_textual_inversion("path-to-save-model", weight_name="<new1>.bin") |
image = pipeline( |
"<new1> cat sitting in a bucket", |
num_inference_steps=100, |
guidance_scale=6.0, |
eta=1.0, |
).images[0] |
image.save("cat.png") |
</hfoption> |
<hfoption id="multiple concepts"> |
Copied import torch |
from huggingface_hub.repocard import RepoCard |
from diffusers import DiffusionPipeline |
pipeline = DiffusionPipeline.from_pretrained("sayakpaul/custom-diffusion-cat-wooden-pot", torch_dtype=torch.float16).to("cuda") |
pipeline.unet.load_attn_procs(model_id, weight_name="pytorch_custom_diffusion_weights.bin") |
pipeline.load_textual_inversion(model_id, weight_name="<new1>.bin") |
pipeline.load_textual_inversion(model_id, weight_name="<new2>.bin") |
image = pipeline( |
"the <new1> cat sculpture in the style of a <new2> wooden pot", |
num_inference_steps=100, |
guidance_scale=6.0, |
eta=1.0, |
).images[0] |
image.save("multi-subject.png") |
</hfoption> |
</hfoptions> |
Next steps Congratulations on training a model with Custom Diffusion! 🎉 To learn more: Read the Multi-Concept Customization of Text-to-Image Diffusion blog post to learn more details about the experimental results from the Custom Diffusion team. |
Text-to-Video Generation with AnimateDiff Overview AnimateDiff: Animate Your Personalized Text-to-Image Diffusion Models without Specific Tuning by Yuwei Guo, Ceyuan Yang, Anyi Rao, Yaohui Wang, Yu Qiao, Dahua Lin, Bo Dai. The abstract of the paper is the following: With the advance of text-to-image models (e.g., Stab... |
from diffusers import AnimateDiffPipeline, DDIMScheduler, MotionAdapter |
from diffusers.utils import export_to_gif |
# Load the motion adapter |
adapter = MotionAdapter.from_pretrained("guoyww/animatediff-motion-adapter-v1-5-2", torch_dtype=torch.float16) |
# load SD 1.5 based finetuned model |
model_id = "SG161222/Realistic_Vision_V5.1_noVAE" |
pipe = AnimateDiffPipeline.from_pretrained(model_id, motion_adapter=adapter, torch_dtype=torch.float16) |
scheduler = DDIMScheduler.from_pretrained( |
model_id, |
subfolder="scheduler", |
clip_sample=False, |
timestep_spacing="linspace", |
beta_schedule="linear", |
steps_offset=1, |
) |
pipe.scheduler = scheduler |
# enable memory savings |
pipe.enable_vae_slicing() |
pipe.enable_model_cpu_offload() |
output = pipe( |
prompt=( |
"masterpiece, bestquality, highlydetailed, ultradetailed, sunset, " |
"orange sky, warm lighting, fishing boats, ocean waves seagulls, " |
"rippling water, wharf, silhouette, serene atmosphere, dusk, evening glow, " |
"golden hour, coastal landscape, seaside scenery" |
), |
negative_prompt="bad quality, worse quality", |
num_frames=16, |
guidance_scale=7.5, |
num_inference_steps=25, |
generator=torch.Generator("cpu").manual_seed(42), |
) |
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