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In the early 1990s, industrial hemp agriculture in North America began with the Hemp Awareness Committee at the University of Manitoba. The Committee worked with the provincial government to get research and development assistance and was able to obtain test plot permits from the Canadian government. Their efforts led to the legalization of industrial hemp (hemp with only minute amounts of tetrahydrocannabinol) in Canada and the first harvest in 1998. In 2017, the cultivated area for hemp in the Prairie provinces include Saskatchewan with more than , Alberta with , and Manitoba with . Canadian hemp is cultivated mostly for its food value as hulled hemp seeds, hemp oils, and hemp protein powders, with only a small fraction devoted to production of hemp fiber used for construction and insulation. France France is Europe's biggest producer (and the world's second largest producer) with cultivated. 70–80% of the hemp fiber produced in 2003 was used for specialty pulp for cigarette papers and technical applications. About 15% was used in the automotive sector, and 5–6% was used for insulation mats. About 95% of hurds were used as animal bedding, while almost 5% was used in the building sector. In 2010–2011, a total of was cultivated with hemp in the EU, a decline compared with previous year. Russia and Ukraine From the 1950s to the 1980s, the Soviet Union was the world's largest producer of hemp ( in 1970). The main production areas were in Ukraine, the Kursk and Orel regions of Russia, and near the Polish border. Since its inception in 1931, the Hemp Breeding Department at the Institute of Bast Crops in Hlukhiv (Glukhov), Ukraine, has been one of the world's largest centers for developing new hemp varieties, focusing on improving fiber quality, per-hectare yields, and low THC content. After the collapse of the Soviet Union, the commercial cultivation of hemp declined sharply. However, at least an estimated 2.5 million acres of hemp grow wild in the Russian Far East and the Black Sea regions.
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United Kingdom In the United Kingdom, cultivation licenses are issued by the Home Office under the Misuse of Drugs Act 1971. When grown for nondrug purposes, hemp is referred to as industrial hemp, and a common product is fiber for use in a wide variety of products, as well as the seed for nutritional aspects and the oil. Feral hemp or ditch weed is usually a naturalized fiber or oilseed strain of Cannabis that has escaped from cultivation and is self-seeding. United States In October 2019, hemp became legal to grow in 46 U.S. states under federal law. As of 2019, 47 states have enacted legislation to make hemp legal to grow at the state level, with several states implementing medical provisions regarding the growing of plants specifically for non-psychoactive CBD. The 2018 Farm Bill, which incorporated the Hemp Farming Act of 2018, removed hemp as a Schedule I drug and instead made it an agricultural commodity. This legalized hemp at the federal level, which made it easier for hemp farmers to get production licenses, acquire loans, and receive federal crop insurance. NH 2014 N.H. Laws, Chap. 18, SD: HB 1008 (2020) S.D. Codified Laws Ann. §38-35-1 et seq. Authorizes the growth, production and transportation of hemp with a license, and directs the Department of Agriculture to submit a state plan to USDA. Requires a minimum of five contiguous outdoor acres for grower license applications, and requires any license applicants to submit to a state and federal criminal background investigation. Requires a transportation permit for any transporter traveling within or through the state and creates two types of industrial hemp transportation permits (grower licensee and general) provided by the Department of Public Safety. Creates the Hemp Regulatory Program Fund.
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The process to legalize hemp cultivation began in 2009, when Oregon began approving licenses for industrial hemp. Then, in 2013, after the legalization of marijuana, several farmers in Colorado planted and harvested several acres of hemp, bringing in the first hemp crop in the United States in over half a century. After that, the federal government created a Hemp Farming Pilot Program as a part of the Agricultural Act of 2014. This program allowed institutions of higher education and state agricultural departments to begin growing hemp without the consent of the Drug Enforcement Administration (DEA). Hemp production in Kentucky, formerly the United States' leading producer, resumed in 2014. Hemp production in North Carolina resumed in 2017, and in Washington State the same year. By the end of 2017, at least 34 U.S. states had industrial hemp programs. In 2018, New York began taking strides in industrial hemp production, along with hemp research pilot programs at Cornell University, Binghamton University and SUNY Morrisville. As of 2017, the hemp industry estimated that annual sales of hemp products were around $820 million annually; hemp-derived CBD have been the major force driving this growth. Despite this progress, hemp businesses in the US have had difficulties expanding as they have faced challenges in traditional marketing and sales approaches. According to a case study done by Forbes, hemp businesses and startups have had difficulty marketing and selling non-psychoactive hemp products, as majority of online advertising platforms and financial institutions do not distinguish between hemp and marijuana. History Gathered hemp fiber was used to make cloth long before agriculture, nine to fifty thousand years ago. It may also be one of the earliest plants to have been cultivated. An archeological site in the Oki Islands of Japan contained cannabis achenes from about 8000 BC, probably signifying use of the plant. Hemp use archaeologically dates back to the Neolithic Age in China, with hemp fiber imprints found on Yangshao culture pottery dating from the 5th millennium BC. The Chinese later used hemp to make clothes, shoes, ropes, and an early form of paper. The classical Greek historian Herodotus (ca. 480 BC) reported that the inhabitants of Scythia would often inhale the vapors of hemp-seed smoke, both as ritual and for their own pleasurable recreation.
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Textile expert Elizabeth Wayland Barber summarizes the historical evidence that Cannabis sativa, "grew and was known in the Neolithic period all across the northern latitudes, from Europe (Germany, Switzerland, Austria, Romania, Ukraine) to East Asia (Tibet and China)," but, "textile use of Cannabis sativa does not surface for certain in the West until relatively late, namely the Iron Age." "I strongly suspect, however, that what catapulted hemp to sudden fame and fortune as a cultigen and caused it to spread rapidly westwards in the first millennium B.C. was the spread of the habit of pot-smoking from somewhere in south-central Asia, where the drug-bearing variety of the plant originally occurred. The linguistic evidence strongly supports this theory, both as to time and direction of spread and as to cause." Jews living in Palestine in the 2nd century were familiar with the cultivation of hemp, as witnessed by a reference to it in the Mishna (Kil'ayim 2:5) as a variety of plant, along with arum, that sometimes takes as many as three years to grow from a seedling. In late medieval Holy Roman Empire (Germany) and Italy, hemp was employed in cooked dishes, as filling in pies and tortes, or boiled in a soup. Hemp in later Europe was mainly cultivated for its fibers and was used for ropes on many ships, including those of Christopher Columbus. The use of hemp as a cloth was centered largely in the countryside, with higher quality textiles being available in the towns.
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The Spaniards brought hemp to the Americas and cultivated it in Chile starting about 1545. Similar attempts were made in Peru, Colombia, and Mexico, but only in Chile did the crop find success. In July 1605, Samuel Champlain reported the use of grass and hemp clothing by the (Wampanoag) people of Cape Cod and the (Nauset) people of Plymouth Bay told him they harvested hemp in their region where it grew wild to a height of 4 to 5 ft. In May 1607, "hempe" was among the crops Gabriel Archer observed being cultivated by the natives at the main Powhatan village, where Richmond, Virginia, is now situated; and in 1613, Samuell Argall reported wild hemp "better than that in England" growing along the shores of the upper Potomac. As early as 1619, the first Virginia House of Burgesses passed an Act requiring all planters in Virginia to sow "both English and Indian" hemp on their plantations. The Puritans are first known to have cultivated hemp in New England in 1645. United States George Washington pushed for the growth of hemp as it was a cash crop commonly used to make rope and fabric. In May 1765 he noted in his diary about the sowing of seeds each day until mid-April. Then he recounts the harvest in October which he grew 27 bushels that year. It is sometimes supposed that an excerpt from Washington's diary, which reads "Began to the Male from the Female hemp at Do.&—rather too late" is evidence that he was trying to grow female plants for the THC found in the flowers. However, the editorial remark accompanying the diary states that "This may arise from their [the male] being coarser, and the stalks larger" In subsequent days, he describes soaking the hemp (to make the fibers usable) and harvesting the seeds, suggesting that he was growing hemp for industrial purposes, not recreational.
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George Washington also imported the Indian hemp plant from Asia, which was used for fiber and, by some growers, for intoxicating resin production. In a 1796 letter to William Pearce who managed the plants for him, Washington says, "What was done with the Indian Hemp plant from last summer? It ought, all of it, to be sown again; that not only a stock of seed sufficient for my own purposes might have been raised, but to have disseminated seed to others; as it is more valuable than common hemp." Other presidents known to have farmed hemp for alternative purposes include Thomas Jefferson, James Madison, James Monroe, Andrew Jackson, Zachary Taylor, and Franklin Pierce. Historically, hemp production had made up a significant portion of antebellum Kentucky's economy. Before the American Civil War, many slaves worked on plantations producing hemp. In 1937, the Marihuana Tax Act of 1937 was passed in the United States, levying a tax on anyone who dealt commercially in cannabis, hemp, or marijuana. The passing of the Act to destroy the U.S. hemp industry has been reputed to involve businessmen Andrew Mellon, Randolph Hearst and the Du Pont family. One claim is that Hearst believed that his extensive timber holdings were threatened by the invention of the decorticator that he feared would allow hemp to become a cheap substitute for the paper pulp used for newspaper. Historical research indicates this fear was unfounded because improvements of the decorticators in the 1930s – machines that separated the fibers from the hemp stem – could not make hemp fiber a cheaper substitute for fibers from other sources. Further, decorticators did not perform satisfactorily in commercial production. Another claim is that Mellon, Secretary of the Treasury and the wealthiest man in America at that time, had invested heavily in DuPont's new synthetic fiber, nylon, and believed that the replacement of the traditional resource, hemp, was integral to the new product's success. DuPont and many industrial historians dispute a link between nylon and hemp, nylon became immediately a scarce commodity. Nylon had characteristics that could be used for toothbrushes (sold from 1938) and very thin nylon fiber could compete with silk and rayon in various textiles normally not produced from hemp fiber, such as very thin stockings for women.
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While the Marijuana Tax Act of 1937 had just been signed into law, the United States Department of Agriculture lifted the tax on hemp cultivation during WWII. Before WWII, the U.S. Navy used Jute and Manila Hemp from the Philippines and Indonesia for the cordage on their ships. During the war, Japan cut off those supply lines. America was forced to turn inward and revitalize the cultivation of Hemp on U.S. soils. Hemp was used extensively by the United States during World War II to make uniforms, canvas, and rope. Much of the hemp used was cultivated in Kentucky and the Midwest. During World War II, the U.S. produced a short 1942 film, Hemp for Victory, promoting hemp as a necessary crop to win the war. By the 1980s the film was largely forgotten, and the U.S. government even denied its existence. The film, and the important historical role of hemp in U.S. agriculture and commerce was brought to light by hemp activist Jack Herer in the book The Emperor Wears No Clothes. U.S. farmers participated in the campaign to increase U.S. hemp production to 36,000 acres in 1942. This increase amounted to more than 20 times the production in 1941 before the war effort. In the United States, Executive Order 12919 (1994) identified hemp as a strategic national product that should be stockpiled. Historical cultivation Hemp has been grown for millennia in Asia and the Middle East for its fiber. Commercial production of hemp in the West took off in the eighteenth century, but was grown in the sixteenth century in eastern England. Because of colonial and naval expansion of the era, economies needed large quantities of hemp for rope and oakum. In the early 1940s, world production of hemp fiber ranged from 250,000 to 350,000 metric tons, Russia was the biggest producer. In Western Europe, the cultivation of hemp was not legally banned by the 1930s, but the commercial cultivation stopped by then, due to decreased demand compared to increasingly popular artificial fibers. Speculation about the potential for commercial cultivation of hemp in large quantities has been criticized due to successful competition from other fibers for many products. The world production of hemp fiber fell from over 300,000 metric tons 1961 to about 75,000 metric tons in the early 1990s and has after that been stable at that level. Japan
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In Japan, hemp was historically used as paper and a fiber crop. There is archaeological evidence cannabis was used for clothing and the seeds were eaten in Japan back to the Jōmon period (10,000 to 300 BC). Many Kimono designs portray hemp, or asa (), as a beautiful plant. In 1948, marijuana was restricted as a narcotic drug. The ban on marijuana imposed by the United States authorities was alien to Japanese culture, as the drug had never been widely used in Japan before. Though these laws against marijuana are some of the world's strictest, allowing five years imprisonment for possession of the drug, they exempt hemp growers, whose crop is used to make robes for Buddhist monks and loincloths for Sumo wrestlers. Because marijuana use in Japan has doubled in the past decade, these exemptions have recently been called into question. Portugal The cultivation of hemp in Portuguese lands began around the fourteenth century. The raw material was used for the preparation of rope and plugs for the Portuguese ships. Portugal also utilized its colonies to support its hemp supply, including in certain parts of Brazil. In order to recover the ailing Portuguese naval fleet after the Restoration of Independence in 1640, King John IV put a renewed emphasis on the growing of hemp. He ordered the creation of the Royal Linen and Hemp Factory in the town of Torre de Moncorvo to increase production and support the effort. In 1971, the cultivation of hemp became illegal, and the production was substantially reduced. Because of EU regulations 1308–70, 619/71 and 1164–89, this law was revoked (for some certified seed varieties).
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Messier 49 (also known as M49 or NGC 4472) is a giant elliptical galaxy about away in the equatorial constellation of Virgo. This galaxy was discovered by astronomer Charles Messier in 1777. As an elliptical galaxy, Messier 49 has the physical form of a radio galaxy, but it only has the radio emission of a normal galaxy. From the detected radio emission, the core region has roughly 1053 erg (1046 J or 1022 YJ) of synchrotron energy. The nucleus of this galaxy is emitting X-rays, suggesting the likely presence of a supermassive black hole with an estimated mass of , or 565 million times the mass of the Sun (). X-ray emissions shows a structure to the north of Messier 49 that resembles a bow shock. To the southwest of the core, the luminous outline of the galaxy can be traced out to a distance of 260 kpc. This galaxy has many globular clusters: estimated to be about 5,900. This is far more than the roughly 200 orbiting the Milky Way, but dwarfed by the 13,450 orbiting the supergiant elliptical galaxy Messier 87. On average, the globular clusters of M49 are about 10 billion years old. Between 2000 and 2009, strong evidence for a stellar mass black hole was discovered in one. A second candidate was announced in 2011. Messier 49 was the first member of the Virgo Cluster of galaxies to be discovered. It is the most luminous member of that cluster and more luminous than any galaxy closer to the Earth. This galaxy forms part of the smaller Virgo B subcluster 4.5° away from the dynamic center of the Virgo Cluster, centered on Messier 87. Messier 49 is gravitationally interacting with the dwarf irregular galaxy UGC 7636. The dwarf shows a trail of debris spanning roughly 1 × 5 arcminutes, which corresponds to a physical dimension of . One supernova has been observed in M49: SN 1969Q (type unknown, mag. 13) was discovered by Evans on 12 June 1969. [Note: some sources incorrectly report the discovery date as 1 June 1969.]
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A fireplace or hearth is a structure made of brick, stone or metal designed to contain a fire. Fireplaces are used for the relaxing ambiance they create and for heating a room. Modern fireplaces vary in heat efficiency, depending on the design. Historically, they were used for heating a dwelling, cooking, and heating water for laundry and domestic uses. A fire is contained in a firebox or fire pit; a chimney or other flue allows exhaust gas to escape. A fireplace may have the following: a foundation, a hearth, a firebox, a mantel, a chimney crane (used in kitchen and laundry fireplaces), a grate, a lintel, a lintel bar, an overmantel, a damper, a smoke chamber, a throat, a flue, and a chimney filter or afterburner. On the exterior, there is often a corbelled brick crown, in which the projecting courses of brick act as a drip course to keep rainwater from running down the exterior walls. A cap, hood, or shroud serves to keep rainwater out of the exterior of the chimney; rain in the chimney is a much greater problem in chimneys lined with impervious flue tiles or metal liners than with the traditional masonry chimney, which soaks up all but the most violent rain. Some chimneys have a spark arrestor incorporated into the crown or cap. Organizations like the United States Environmental Protection Agency (EPA) and the Washington State Department of Ecology warn that, according to various studies, fireplaces can pose health risks. The EPA writes "Smoke may smell good, but it's not good for you." Types of fireplaces Manufactured fireplaces are made with sheet metal or glass fire boxes. Electric fireplaces can be built-in replacements for wood or gas or retrofit with log inserts or electric fireboxes. A few types are wall mounted electric fireplaces, electric fireplace stoves, electric mantel fireplaces, and fixed or free standing electric fireplaces. Masonry and prefabricated fireplaces can be fueled by: Wood fuel or firewood and other biomass Charcoal (carbonized biomass) Coal of various grades Coke (carbonized coal) Smokeless fuel of several types Flammable gases: propane, butane, and methane (natural gas is mostly methane, liquefied petroleum gas mostly propane) Ethanol (a liquid alcohol, also sold in gels)
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Ventless fireplaces (duct free/room-venting fireplaces) are fueled by either gel, liquid propane, bottled gas or natural gas. In the United States, some states and local counties have laws restricting these types of fireplaces. They must be properly sized to the area to be heated. There are also air quality control issues due to the amount of moisture they release into the room air, and an oxygen sensor and a carbon monoxide detector are safety essentials. Direct vent fireplaces are fueled by either liquid propane or natural gas. They are completely sealed from the area that is heated, and vent all exhaust gasses to the exterior of the structure. Chimney and flue types: Masonry (brick or stone fireplaces and chimneys) with or without tile-lined flue. Reinforced concrete chimneys. Fundamental design flaws bankrupted the US manufacturers and made the design obsolete. These chimneys often show vertical cracks on the exterior. Metal-lined flue: Double- or triple-walled metal pipe running up inside a new or existing wood-framed or masonry chase. Newly constructed flues may feature a chase cover, a cap, and a spark arrestor at the top to keep small animals out and to prevent sparks from being broadcast into the atmosphere. All gas fireplaces require trained gas service members to carry out installations. Accessories A wide range of accessories are used with fireplaces, which range between countries, regions, and historical periods. For the interior, common in recent Western cultures include grates, fireguards, log boxes, andirons and pellet baskets, all of which cradle fuel and accelerate combustion. A grate (or fire grate) is a frame, usually of iron bars, to retain fuel for a fire. Heavy metal firebacks are sometimes used to capture and re-radiate heat, to protect the back of the fireplace, and as decoration. Fenders are low metal frames set in front of the fireplace to contain embers, soot and ash. For fireplace tending, tools include pokers, bellows, tongs, shovels, brushes and tool stands. Other wider accessories can include log baskets, companion sets, coal buckets, cabinet accessories and more. History Ancient fire pits were sometimes built in the ground, within caves, or in the center of a hut or dwelling. Evidence of prehistoric, man-made fires exists on all six inhabited continents. The disadvantage of early indoor fire pits was that they produced toxic and/or irritating smoke inside the dwelling.
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Fire pits developed into raised hearths in buildings, but venting smoke depended on open windows or holes in roofs. The medieval great hall typically had a centrally located hearth, where an open fire burned with the smoke rising to the vent in the roof. Louvers were developed during the Middle Ages to allow the roof vents to be covered so rain and snow would not enter. Also during the Middle Ages, smoke canopies were invented to prevent smoke from spreading through a room and vent it out through a wall or roof. These could be placed against stone walls, instead of taking up the middle of the room, and this allowed smaller rooms to be heated. Chimneys were invented in northern Europe in the 11th or 12th century and largely fixed the problem of smoke, more reliably venting it outside. They made it possible to give the fireplace a draft, and also made it possible to put fireplaces in multiple rooms in buildings conveniently. They did not come into general use immediately, however, as they were expensive to build and maintain. In 1678, Prince Rupert, nephew of Charles I, raised the grate of the fireplace, improving the airflow and venting system. The 18th century saw two important developments in the history of fireplaces. Benjamin Franklin developed a convection chamber for the fireplace that greatly improved the efficiency of fireplaces and wood-burning stoves. He also improved the airflow by pulling air from a basement and venting out a longer area at the top. In the later 18th century, Count Rumford designed a fireplace with a tall, shallow firebox that was better at drawing the smoke up and out of the building. The shallow design also improved greatly the amount of heat transfer projected into the room. The Aesthetic movement of the 1870s and 1880s favoured a more traditional look based on stone, with simple designs and limited ornamentation. In the 1890s, the Aesthetic movement gave way to the Arts and Crafts movement, which still emphasized quality stone and practical features. Stone fireplaces at this time were a symbol of prosperity, as to some degree they remain today. Evolution of fireplace design Over time, the purpose of fireplaces has changed from one of necessity to one of visual interest. Early ones were more fire pits than modern fireplaces. They were used for warmth on cold days and nights, as well as for cooking. They also served as a gathering place within the home. These fire pits were usually centered within a room, allowing more people to gather around it.
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Many flaws were found in early fireplace designs. Along with the Industrial Revolution, came large-scale housing developments, necessitating a standardization of fireplaces. The most renowned fireplace designers of this time were the Adam Brothers: John Adam, Robert Adam, and James Adam. They perfected a style of fireplace design that was used for generations. It was smaller, more brightly lit, with an emphasis on the quality of the materials used in their construction, instead of their size. By the 1800s, most new fireplaces were made up of two parts, the surround and the insert. The surround consisted of the mantelpiece and side supports, usually in wood, marble or granite. The insert was where the fire burned, and was constructed of cast iron often backed with decorative tiles. As well as providing heat, the fireplaces of the Victorian era were thought to add a cosy ambiance to homes. In the US state of Wisconsin, some elementary classrooms would contain decorated fireplaces to ease children's transition from home to school. Heating efficiency Some fireplace units incorporate a blower, which transfers more of the fireplace's heat to the air via convection, resulting in a more evenly heated space and a lower heating load. Fireplace efficiency can also be increased with the use of a fireback, a piece of metal that sits behind the fire and reflects heat back into the room. Firebacks are traditionally made from cast iron, but are also made from stainless steel. Most older fireplaces have a relatively low efficiency rating. Standard, modern, wood-burning masonry fireplaces though have an efficiency rating of at least 80% (legal minimum requirement, for example, in Salzburg, Austria). To improve efficiency, fireplaces can also be modified by inserting special heavy fireboxes designed to burn much cleaner and can reach efficiencies as high as 80% in heating the air. These modified fireplaces are often equipped with a large fire window, enabling an efficient heating process in two phases. During the first phase the initial heat is provided through a large glass window while the fire is burning. During this time the structure, built of refractory bricks, absorbs the heat. This heat is then evenly radiated for many hours during the second phase. Masonry fireplaces without a glass fire window only provide heat radiated from its surface. Depending on the outside temperature, 1 to 2 daily firings are sufficient to ensure a constant room temperature. Health effects
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Wood A literature review published in the Journal of Toxicology and Environmental Health concludes that there are a wide variety of health risks posed by residential wood combustion. It states: The Washington State Department of Ecology also published a booklet explaining why wood smoke can be dangerous. It explains that human lung and respiratory systems are unable to filter particulates emitted by wood combustion, which penetrate deeply into the lungs. For months, carcinogens can continue to cause changes and structural damage within the respiratory system. Young children, seniors, pregnant women, smokers and individuals with respiratory diseases are most vulnerable. Wood smoke can cause disease and even death in children, because it is associated with lower respiratory tract infections. Home fireplaces have caused fatal carbon monoxide poisoning. Gases and ethanol Propane, butane, and methane are all flammable gases used in fireplaces (natural gas is mostly methane, liquefied petroleum gas mostly propane). Gases can act as asphyxiant gases or cause gas explosions if they are allowed to accumulate unburned. Ethanol (a liquid, also sold in gels) fires can also cause severe burns. Burning hydrocarbons can decrease indoor air quality. Emissions include airborne particulate matter (such as black carbon) and gases like nitrogen oxide. These harm health: they weaken the immune system, and increase infections, blood pressure, cardiovascular diseases, and insulin resistance. Some forms of fuel are more harmful than others. Burning hydrocarbon fuels incompletely can produce carbon monoxide, which is highly poisonous and can cause death and long-term neurological disorders. Environmental effects Burning any hydrocarbon fuel releases carbon dioxide and water vapor. Other emissions, such as nitrogen oxides and sulfur oxides, can be harmful to the environment. Glossary
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Several of these terms may be compounded with chimney or fireplace such as chimney-back. Andiron—Either one of two horizontal metal bars resting on short legs intended to support firewood in a hearth. Arch—An arched top of the fireplace opening. Ash dump—An opening in a hearth to sweep ashes for later removal from the ash pit. Back (fireback)—The inside, rear wall of the fireplace of masonry or metal that reflects heat into the room. Brick trimmer—A brick arch supporting a hearth or shielding a joist in front of a fireplace. Chimney breast—The part of the chimney which projects into a room to accommodate a fireplace. Crane—Metal arms mounted on pintles, which swing and hold pots above a fire. Damper—A metal door to close a flue when a fireplace is not in use. Flue—The passageway in the chimney. Hearth—The floor of a fireplace. The part of a hearth which projects into a room may be called the front or outer hearth. Hearthstone—A large stone or other materials used as the hearth material. Insert—The fireplace insert is a device inserted into an existing masonry or prefabricated wood fireplace. Jamb—The side of a fireplace opening. Mantel—Either the shelf above a fireplace or the structure to support masonry above a fireplace Smoke shelf—A shelf below the smoke chamber and behind the damper. It collects debris and water falling down the flue. Throat (waist)—The narrow area above a fireplace usually where the damper is located. Wing—The sides of a fireplace above the opening near the throat.
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A scale or balance is a device used to measure weight or mass. These are also known as mass scales, weight scales, mass balances, massometers, and weight balances. The traditional scale consists of two plates or bowls suspended at equal distances from a fulcrum. One plate holds an object of unknown mass (or weight), while objects of known mass or weight, called weights, are added to the other plate until mechanical equilibrium is achieved and the plates level off, which happens when the masses on the two plates are equal. The perfect scale rests at neutral. A spring scale will make use of a spring of known stiffness to determine mass (or weight). Suspending a certain mass will extend the spring by a certain amount depending on the spring's stiffness (or spring constant). The heavier the object, the more the spring stretches, as described in Hooke's law. Other types of scales making use of different physical principles also exist. Some scales can be calibrated to read in units of force (weight) such as newtons instead of units of mass such as kilograms. Scales and balances are widely used in commerce, as many products are sold and packaged by mass. Pan balance History The balance scale is such a simple device that its usage likely far predates the evidence. What has allowed archaeologists to link artifacts to weighing scales are the stones for determining absolute mass. The balance scale itself was probably used to determine relative mass long before absolute mass. The oldest attested evidence for the existence of weighing scales dates to the Fourth Dynasty of Egypt, with Deben (unit) balance weights, from the reign of Sneferu (c. 2600 BC) excavated, though earlier usage has been proposed. Carved stones bearing marks denoting mass and the Egyptian hieroglyphic symbol for gold have been discovered, which suggests that Egyptian merchants had been using an established system of mass measurement to catalog gold shipments or gold mine yields. Although no actual scales from this era have survived, many sets of weighing stones as well as murals depicting the use of balance scales suggest widespread usage.
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Examples, dating , have also been found in the Indus River valley. Uniform, polished stone cubes discovered in early settlements were probably used as mass-setting stones in balance scales. Although the cubes bear no markings, their masses are multiples of a common denominator. The cubes are made of many different kinds of stones with varying densities. Clearly their mass, not their size or other characteristics, was a factor in sculpting these cubes. In China, the earliest weighing balance excavated was from a tomb of the State of Chu of the Chinese Warring States Period dating back to the 3rd to 4th century BC in Mount Zuojiagong near Changsha, Hunan. The balance was made of wood and used bronze masses. Variations on the balance scale, including devices like the cheap and inaccurate bismar (unequal-armed scales), began to see common usage by c. 400 BC by many small merchants and their customers. A plethora of scale varieties each boasting advantages and improvements over one another appear throughout recorded history, with such great inventors as Leonardo da Vinci lending a personal hand in their development. Even with all the advances in weighing scale design and development, all scales until the seventeenth century AD were variations on the balance scale. The standardization of the weights used – and ensuring traders used the correct weights – was a considerable preoccupation of governments throughout this time. The original form of a balance consisted of a beam with a fulcrum at its center. For highest accuracy, the fulcrum would consist of a sharp V-shaped pivot seated in a shallower V-shaped bearing. To determine the mass of the object, a combination of reference masses was hung on one end of the beam while the object of unknown mass was hung on the other end (see balance and steelyard balance). For high precision work, such as empirical chemistry, the center beam balance is still one of the most accurate technologies available, and is commonly used for calibrating test masses. However, bronze fragments discovered in central Germany and Italy had been used during the Bronze Age as an early form of currency. In the same time period, merchants had used standard weights of equivalent value between 8 and 10.5 grams from Great Britain to Mesopotamia. Mechanical balances
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The balance (also balance scale, beam balance and laboratory balance) was the first mass measuring instrument invented. In its traditional form, it consists of a pivoted horizontal lever with arms of equal lengththe beam or tron and a weighing pan suspended from each arm (hence the plural name "scales for a weighing instrument). The unknown mass is placed in one pan and standard masses are added to the other pan until the beam is as close to equilibrium as possible. In precision balances, a more accurate determination of the mass is given by the position of a sliding mass moved along a graduated scale. A decimal balance uses the lever in which the arm for weights is 10 times longer than the arm for weighted objects, so that much lighter weights may be used to weigh heavy object. Similarly a centesimal balance uses arms in ratio 1:100. Unlike spring-based scales, balances are used for the precision measurement of mass as their accuracy is not affected by variations in the local gravitational field. (On Earth, for example, these can amount to ±0.5% between locations.) A change in the strength of the gravitational field caused by moving the balance does not change the measured mass, because the moments of force on either side of the center balanced beam are affected equally. A center beam balance will render an accurate measurement of mass at any location experiencing a constant gravity or acceleration. Very precise measurements are achieved by ensuring that the balance's fulcrum is essentially friction-free (a knife edge is the traditional solution), by attaching a pointer to the beam which amplifies any deviation from a balance position; and finally by using the lever principle, which allows fractional masses to be applied by movement of a small mass along the measuring arm of the beam, as described above. For greatest accuracy, there needs to be an allowance for the buoyancy in air, whose effect depends on the densities of the masses involved.
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To reduce the need for large reference masses, an off-center beam can be used. A balance with an off-center beam can be almost as accurate as a scale with a center beam, but the off-center beam requires special reference masses and cannot be intrinsically checked for accuracy by simply swapping the contents of the pans as a center-beam balance can. To reduce the need for small graduated reference masses, a sliding weight called a poise can be installed so that it can be positioned along a calibrated scale. A poise adds further intricacies to the calibration procedure, since the exact mass of the poise must be adjusted to the exact lever ratio of the beam. For greater convenience in placing large and awkward loads, a platform can be floated on a cantilever beam system which brings the proportional force to a noseiron bearing; this pulls on a stilyard rod to transmit the reduced force to a conveniently sized beam. One still sees this design in portable beam balances of 500 kg capacity which are commonly used in harsh environments without electricity, as well as in the lighter duty mechanical bathroom scale (which actually uses a spring scale, internally). The additional pivots and bearings all reduce the accuracy and complicate calibration; the float system must be corrected for corner errors before the span is corrected by adjusting the balance beam and poise. Roberval balance In 1669 the Frenchman Gilles Personne de Roberval presented a new kind of balance scale to the French Academy of Sciences. This scale consisted of a pair of vertical columns separated by a pair of equal-length arms and pivoting in the center of each arm from a central vertical column, creating a parallelogram. From the side of each vertical column a peg extended. To the amazement of observers, no matter where Roberval hung two equal weight along the peg, the scale still balanced. In this sense, the scale was revolutionary: it evolved into the more-commonly encountered form consisting of two pans placed on vertical column located above the fulcrum and the parallelogram below them. The advantage of the Roberval design is that no matter where equal weights are placed in the pans, the scale will still balance.
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Further developments have included a "gear balance" in which the parallelogram is replaced by any odd number of interlocking gears greater than one, with alternating gears of the same size and with the central gear fixed to a stand and the outside gears fixed to pans, as well as the "sprocket gear balance" consisting of a bicycle-type chain looped around an odd number of sprockets with the central one fixed and the outermost two free to pivot and attached to a pan. Because it has more moving joints which add friction, the Roberval balance is consistently less accurate than the traditional beam balance, but for many purposes this is compensated for by its usability. Torsion balance The torsion balance is one of the most mechanically accurate of analog balances. Pharmacy schools still teach how to use torsion balances in the U.S. It utilizes pans like a traditional balance that lie on top of a mechanical chamber which bases measurements on the amount of twisting of a wire or fiber inside the chamber. The scale must still use a calibration weight to compare against, and can weigh objects greater than 120 mg and come within a margin of error +/- 7 mg. Many microbalances and ultra-microbalances that weigh fractional gram values are torsion balances. A common fiber type is quartz crystal. Electronic devices Microbalance A microbalance (also called an ultramicrobalance, or nanobalance) is an instrument capable of making precise measurements of the mass of objects of relatively small mass: on the order of a million parts of a gram and below. Analytical balance
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An analytical balance is a class of balance designed to measure small mass in the sub-milligram range. The measuring pan of an analytical balance (0.1 mg or better) is inside a transparent enclosure with doors so that dust does not collect and so any air currents in the room do not affect the balance's operation. This enclosure is often called a draft shield. The use of a mechanically vented balance safety enclosure, which has uniquely designed acrylic airfoils, allows a smooth turbulence-free airflow that prevents balance fluctuation and the measure of mass down to 1 μg without fluctuations or loss of product. Also, the sample must be at room temperature to prevent natural convection from forming air currents inside the enclosure from causing an error in reading. Single-pan mechanical substitution balances maintain consistent response throughout the useful capacity, which is achieved by maintaining a constant load on the balance beam and thus the fulcrum by subtracting mass on the same side of the beam to which the sample is added. Electronic analytical scales measure the force needed to counter the mass being measured rather than using actual masses. As such they must have calibration adjustments made to compensate for gravitational differences. They use an electromagnet to generate a force to counter the sample being measured and output the result by measuring the force needed to achieve balance. Such a measurement device is called an electromagnetic force restoration sensor. Pendulum balance scales Pendulum type scales do not use springs. These designs use pendulums and operate as a balance that is unaffected by differences in gravity. An example of application of this design are scales made by the Toledo Scale Company. Programmable scales A programmable scale has a programmable logic controller in it, allowing it to be programmed for various applications such as batching, labeling, filling (with check weight function), truck scales, and more. Another important function is counting, e. g. used to count small parts in larger quantities during the annual stock taking. Counting scales (which can also do just weighing) can range from mg to tonnes. Symbolism The scales (specifically, a two-pan, beam balance) are one of the traditional symbols of justice, as wielded by statues of Lady Justice. This corresponds to the use in a metaphor of matters being "held in the balance". It has its origins in ancient Egypt.
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Scales also are widely used as a symbol of finance, commerce, or trade, in which they have played a traditional, vital role since ancient times. For instance, balance scales are depicted in the seal of the U.S. Department of the Treasury and the Federal Trade Commission. Scales are also the symbol for the astrological sign Libra. Scales (specifically, a two-pan, beam balance in a state of equal balance) are the traditional symbol of Pyrrhonism indicating the equal balance of arguments used in inducing epoche. Force-measuring (weight) scales History Although records dating to the 1700s refer to spring scales for measuring mass, the earliest design for such a device dates to 1770 and credits Richard Salter, an early scale-maker. Spring scales came into wide usage in the United Kingdom after 1840 when R. W. Winfield developed the candlestick scale for weighing letters and packages, required after the introduction of the Uniform Penny Post. Postal workers could work more quickly with spring scales than balance scales because they could be read instantaneously and did not have to be carefully balanced with each measurement. By the 1940s, various electronic devices were being attached to these designs to make readings more accurate. Load cells – transducers that convert force to an electrical signal – have their beginnings as early as the late nineteenth century, but it was not until the late twentieth century that their widespread usage became economically and technologically viable. Mechanical scales A mechanical scale or balance is used to describe a weighing device that is used to measure the mass, force exertion, tension, and resistance of an object without the need of a power supply. Types of mechanical scales include decimal balances, spring scales, hanging scales, triple beam balances, and force gauges. Spring scales A spring scale measures mass by reporting the distance that a spring deflects under a load. This contrasts to a balance, which compares the torque on the arm due to a sample weight to the torque on the arm due to a standard reference mass using a horizontal lever. Spring scales measure force, which is the tension force of constraint acting on an object, opposing the local force of gravity. They are usually calibrated so that measured force translates to mass at earth's gravity. The object to be weighed can be simply hung from the spring or set on a pivot and bearing platform.
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In a spring scale, the spring either stretches (as in a hanging scale in the produce department of a grocery store) or compresses (as in a simple bathroom scale). By Hooke's law, every spring has a proportionality constant that relates how hard it is pulled to how far it stretches. Weighing scales use a spring with a known spring constant (see Hooke's law) and measure the displacement of the spring by any variety of mechanisms to produce an estimate of the gravitational force applied by the object. Rack and pinion mechanisms are often used to convert the linear spring motion to a dial reading. Spring scales have two sources of error that balances do not: the measured mass varies with the strength of the local gravitational force (by as much as 0.5% at different locations on Earth), and the elasticity of the measurement spring can vary slightly with temperature. With proper manufacturing and setup, however, spring scales can be rated as legal for commerce. To remove the temperature error, a commerce-legal spring scale must either have temperature-compensated springs or be used at a fairly constant temperature. To eliminate the effect of gravity variations, a commerce-legal spring scale must be calibrated where it is used. Hydraulic or pneumatic scale It is also common in high-capacity applications such as crane scales to use hydraulic force to sense mass. The test force is applied to a piston or diaphragm and transmitted through hydraulic lines to a dial indicator based on a Bourdon tube or electronic sensor. Domestic Weighing Scale Electronic digital scales display weight as a number, usually on a liquid crystal display (LCD). They are versatile because they may perform calculations on the measurement and transmit it to other digital devices. On a digital scale, the force of the weight causes a spring to deform, and the amount of deformation is measured by one or more transducers called strain gauges. A strain gauge is a conductor whose electrical resistance changes when its length changes. Strain gauges have limited capacity and larger digital scales may use a hydraulic transducer called a load cell instead. A voltage is applied to the device, and the weight causes the current through it to change. The current is converted to a digital number by an analog-to-digital converter, translated by digital logic to the correct units, and displayed on the display. Usually, the device is run by a microprocessor chip.
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Digital bathroom scale A digital bathroom scale is a scale on the floor which a person stands on. The weight is shown on an LED or LCD display. The digital electronics may do more than just display weight, it may calculate body fat, BMI, lean mass, muscle mass, and water ratio. Some modern bathroom scales are wirelessly or cellularly connected and have features like smartphone integration, cloud storage, and fitness tracking. They are usually powered by a button cell, or battery of AA or AAA size. Digital kitchen scale Digital kitchen scales are used for weighing food in a kitchen during cooking. These are usually lightweight and compact. Strain gauge scale In electronic versions of spring scales, the deflection of a beam supporting the unknown mass is measured using a strain gauge, which is a length-sensitive electrical resistance. The capacity of such devices is only limited by the resistance of the beam to deflection. The results from several supporting locations may be added electronically, so this technique is suitable for determining the mass of very heavy objects, such as trucks and rail cars, and is used in a modern weighbridge. Supermarket and other retail scale These scales are used in the modern bakery, grocery, delicatessen, seafood, meat, produce and other perishable goods departments. Supermarket scales can print labels and receipts, mark mass and count, unit price, total price and in some cases tare. Some modern supermarket scales print an RFID tag that can be used to track the item for tampering or returns. In most cases, these types of scales have a sealed calibration so that the reading on the display is correct and cannot be tampered with. In the US, the scales are certified by the National Type Evaluation Program (NTEP), in South Africa by the South African Bureau of Standards, in Australia, they are certified by the National Measurement Institute (NMI) and in the UK by the International Organization of Legal Metrology. Industrial weighing scale An industrial weighing scale is a device that measures the weight or mass of objects in various industries. It can range from small bench scales to large weighbridges, and it can have different features and capacities. Industrial weighing scales are used for quality control, inventory management, and trade purposes. There are many kinds of industrial weighing scales that are used for different purposes and applications. Some of the common types are:
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Weighbridges : A large scale that can weigh trucks, lorries, containers, and other heavy-duty vehicles. They are used in industries like manufacturing, shipping, mining, agriculture, etc Container Stacker Scale : A container stacker scale is a specialized weighing system designed for accurately measuring the weight of shipping containers. It is typically integrated into the equipment used for loading and unloading containers, such as container handlers or stacker cranes. Container stacker scales provide real-time weight measurements, allowing logistics professionals to ensure that each container is loaded within the specified weight limits. Container stacker scales are used in industries like ports, shipping, and logistics Forklift scale : A forklift scale is a weighing system that is built into a forklift truck. It allows for the weighing of loads while they are being lifted and transported by the forklift. This eliminates the need for separate weighing operations and reduces the time and labor required for material handling operations. Forklift scales are used in various industries, such as manufacturing, logistics, and shipping. Material Handler Scale : A Material Handler Scale is a weighing system that is integrated into a material handler machine, such as a grapple or a magnet. It allows for the accurate and efficient weighing of materials while they are being moved, unloaded, or loaded. A Material Handler Scale can be used in various industries, such as scrap, recycling, waste, and port and harbor. A Material Handler Scale can also transfer the weighing information to a cloud service or an ERP system for real-time monitoring and management of material flow. A pallet jack scale is a device that combines a pallet jack and a weighing scale. It allows you to weigh and move pallets at the same time, saving time and labor. Pallet jack scales are used in various industries, such as manufacturing, logistics, and shipping. Crane Scale : A crane scale is a device that measures the weight or mass of objects that are suspended from a crane. It has a hook at the bottom and a large display that allows distant viewing. Crane scales are used for various industrial applications, such as manufacturing, shipping, mining, recycling, and more
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Wheel Loader Scale : A wheel loader scale is a system that measures the weight of the materials lifted by a wheel loader, a type of heavy machinery used for moving large amounts of earth, sand, gravel, or other materials. A wheel loader scale can help improve the efficiency and accuracy of loading operations, as well as the inventory management and safety of the industries that use them. A wheel loader scale typically consists of a hydraulic sensor, a display unit, and a data management system. The hydraulic sensor is installed in the wheel loader and detects the pressure changes caused by the load. The display unit shows the weight information to the operator and allows them to set target loads, select products and customers, and export data. The data management system can store, analyze, and transmit the weight data to other devices or platforms. Testing and certification Most countries regulate the design and servicing of scales used for commerce. For example, in the European Union weighing instruments are subject to 2014/31/EU and 2014/32/EU directives. A conformity assessment procedure is carried out before placing the instrument on the market, andv the instruments are verified after a given period of time in member states of the European Union. This has tended to cause scale technology to lag behind other technologies because expensive regulatory hurdles are involved in introducing new designs. Nevertheless, there has been a trend to "digital load cells" which are actually strain-gauge cells with dedicated analog converters and networking built into the cell itself. Such designs have reduced the service problems inherent with combining and transmitting a number of 20 millivolt signals in hostile environments. Government regulation generally requires periodic inspections by licensed technicians, using masses whose calibration is traceable to an approved laboratory. Scales intended for non-trade use, such as those used in bathrooms, doctor's offices, kitchens (portion control), and price estimation (but not official price determination) may be produced, but must by law be labelled "Not Legal for Trade" to ensure that they are not re-purposed in a way that jeopardizes commercial interest. In the United States, the document describing how scales must be designed, installed, and used for commercial purposes is NIST Handbook 44. Legal For Trade (LFT) certification usually approve the readability by testing repeatability of measurements to ensure a maximum margin of error of 10%.
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Because gravity varies by over 0.5% over the surface of the earth, the distinction between force due to gravity and mass is relevant for accurate calibration of scales for commercial purposes. Usually, the goal is to measure the mass of the sample rather than its force due to gravity at that particular location. Traditional mechanical balance-beam scales intrinsically measured mass. But ordinary electronic scales intrinsically measure the gravitational force between the sample and the earth, i.e. the weight of the sample, which varies with location. So such a scale has to be re-calibrated after installation, for that specific location, in order to obtain an accurate indication of mass. Sources of error Some of the sources of error in weighing are: Buoyancy – Objects in air develop a buoyancy force that is directly proportional to the volume of air displaced. The difference in density of air due to barometric pressure and temperature creates errors. Error in the mass of reference weight Air gusts, even small ones, which push the scale up or down Friction in the moving components that causes the scale to reach equilibrium at a different configuration than a frictionless equilibrium should occur. Settling airborne dust contributing to the weight Mis-calibration over time, due to drift in the circuit's accuracy, or temperature change Mis-aligned mechanical components due to thermal expansion or contraction of components Magnetic fields acting on ferrous components Forces from electrostatic fields, for example, from feet shuffled on carpets on a dry day Chemical reactivity between air and the substance being weighed (or the balance itself, in the form of corrosion) Condensation of atmospheric water on cold items Evaporation of water from wet items Convection of air from hot or cold items Gravitational differences for a scale which measures force, but not for a balance. Vibration and seismic disturbances Hybrid spring and balance scales
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Elastic arm scale In 2014 a concept of hybrid scale was introduced, the elastically deformable arm scale, which is a combination between a spring scale and a beam balance, exploiting simultaneously both principles of equilibrium and deformation. In this scale, the rigid arms of a classical beam balance (for example a steelyard) are replaced with a flexible elastic rod in an inclined frictionless sliding sleeve. The rod can reach a unique sliding equilibrium when two vertical dead loads (or masses) are applied at its edges. Equilibrium, which would be impossible with rigid arms, is guaranteed because configurational forces develop at the two edges of the sleeve as a consequence of both the free sliding condition and the nonlinear kinematics of the elastic rod. This mass measuring device can also work without a counterweight.
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A dwarf galaxy is a small galaxy composed of about 1000 up to several billion stars, as compared to the Milky Way's 200–400 billion stars. The Large Magellanic Cloud, which closely orbits the Milky Way and contains over 30 billion stars, is sometimes classified as a dwarf galaxy; others consider it a full-fledged galaxy. Dwarf galaxies' formation and activity are thought to be heavily influenced by interactions with larger galaxies. Astronomers identify numerous types of dwarf galaxies, based on their shape and composition. Formation One theory states that most galaxies, including dwarf galaxies, form in association with dark matter, or from gas that contains metals. However, NASA's Galaxy Evolution Explorer space probe identified new dwarf galaxies forming out of gases with low metallicity. These galaxies were located in the Leo Ring, a cloud of hydrogen and helium around two massive galaxies in the constellation Leo. Because of their small size, dwarf galaxies have been observed being pulled toward and ripped by neighbouring spiral galaxies, resulting in stellar streams and eventually galaxy merger. Local dwarf galaxies There are many dwarf galaxies in the Local Group; these small galaxies frequently orbit larger galaxies, such as the Milky Way, the Andromeda Galaxy and the Triangulum Galaxy. A 2007 paper has suggested that many dwarf galaxies were created by galactic tides during the early evolutions of the Milky Way and Andromeda. Tidal dwarf galaxies are produced when galaxies collide and their gravitational masses interact. Streams of galactic material are pulled away from the parent galaxies and the halos of dark matter that surround them. A 2018 study suggests that some local dwarf galaxies formed extremely early, during the Dark Ages within the first billion years after the Big Bang. More than 20 known dwarf galaxies orbit the Milky Way, and recent observations have also led astronomers to believe the largest globular cluster in the Milky Way, Omega Centauri, is in fact the core of a dwarf galaxy with a black hole at its centre, which was at some time absorbed by the Milky Way. Common types Elliptical galaxy: dwarf elliptical galaxy (dE) Dwarf spheroidal galaxy (dSph): Once a subtype of dwarf ellipticals, now regarded as a distinct type Irregular galaxy: dwarf irregular galaxy (dIrr) Spiral galaxy: dwarf spiral galaxy (dS) Magellanic type dwarfs Blue compact dwarf galaxies (see section below) Ultra-compact dwarf galaxies (see section below) Blue compact dwarf galaxies
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In astronomy, a blue compact dwarf galaxy (BCD galaxy) is a small galaxy which contains large clusters of young, hot, massive stars. These stars, the brightest of which are blue, cause the galaxy itself to appear blue in colour. Most BCD galaxies are also classified as dwarf irregular galaxies or as dwarf lenticular galaxies. Because they are composed of star clusters, BCD galaxies lack a uniform shape. They consume gas intensely, which causes their stars to become very violent when forming. BCD galaxies cool in the process of forming new stars. The galaxies' stars are all formed at different time periods, so the galaxies have time to cool and to build up matter to form new stars. As time passes, this star formation changes the shape of the galaxies. Nearby examples include NGC 1705, NGC 2915, NGC 3353 and UGCA 281. Ultra-faint dwarf galaxies Ultra-faint dwarf galaxies (UFDs) are a class of galaxies that contain from a few hundred to one hundred thousand stars, making them the faintest galaxies in the Universe. UFDs resemble globular clusters (GCs) in appearance but have very different properties. Unlike GCs, UFDs contain a significant amount of dark matter and are more extended. UFDs were first discovered with the advent of digital sky surveys in 2005, in particular with the Sloan Digital Sky Survey (SDSS). UFDs are the most dark matter-dominated systems known. Astronomers believe that UFDs encode valuable information about the early Universe, as all UFDs discovered so far are ancient systems that have likely formed very early on, only a few million years after the Big Bang and before the epoch of reionization. Recent theoretical work has hypothesised the existence of a population of young UFDs that form at a much later time than the ancient UFDs. These galaxies have not been observed in our Universe so far.
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Ultra-compact dwarfs Ultra-compact dwarf galaxies (UCD) are a class of very compact galaxies with very high stellar densities, discovered in the 2000s. They are thought to be on the order of 200 light years across, containing about 100 million stars. It is theorised that these are the cores of nucleated dwarf elliptical galaxies that have been stripped of gas and outlying stars by tidal interactions, travelling through the hearts of rich clusters. UCDs have been found in the Virgo Cluster, Fornax Cluster, Abell 1689, and the Coma Cluster, amongst others. In particular, an unprecedentedly large sample of ~ 100 UCDs has been found in the core region of the Virgo cluster by the Next Generation Virgo Cluster Survey team. The first ever relatively robust studies of the global properties of Virgo UCDs suggest that UCDs have distinct dynamical and structural properties from normal globular clusters. An extreme example of UCD is M60-UCD1, about 54 million light years away, which contains approximately 200 million solar masses within a 160 light year radius; the stars in its central region are packed 25 times more densely than stars in Earth's region in the Milky Way. M59-UCD3 is approximately the same size as M60-UCD1 with a half-light radius, rh, of approximately 20 parsecs but is 40% more luminous with an absolute visual magnitude of approximately −14.6. This makes M59-UCD3 the second densest known galaxy. Based on stellar orbital velocities, two UCD in the Virgo Cluster are claimed to have supermassive black holes weighing 13% and 18% of the galaxies' masses. Partial list Aquarius Dwarf Canis Major Dwarf Galaxy Carina Dwarf Crater 2 dwarf Draco Dwarf Eridanus II Fornax Dwarf Henize 2-10 I Zwicky 18 IC 10 Large Magellanic Cloud Leo I Leo II NGC 1569 NGC 1705 NGC 2915 NGC 3353 Pegasus Dwarf Irregular Galaxy PHL 293B Phoenix Dwarf Sagittarius Dwarf Spheroidal Galaxy Sagittarius Dwarf Irregular Galaxy Sculptor Dwarf Galaxy Sculptor Dwarf Irregular Galaxy Sextans A Sextans Dwarf Spheroidal Small Magellanic Cloud Tucana Dwarf Ursa Major I Dwarf Ursa Major II Dwarf Ursa Minor Dwarf Willman 1 Gallery
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The Spanish dancer, scientific name Hexabranchus sanguineus (literally meaning "blood-colored six-gills"), is a dorid nudibranch, a very large and colorful sea slug, a marine gastropod mollusk in the family Hexabranchidae.The taxonomy of the genus Hexabranchus has been controversial but a thorough molecular and morphological study published in 2023 showed that the name H. sanguineus was being used for at least 5 distinct species. Description Hexabranchus sanguineus is a large dorid nudibranch which commonly grows up to a maximum length of 25 cm, with some reports to 40 cm in the Red Sea. All Hexabranchus species have soft, flattened bodies, the anterior dorsal portion has a pair of retractable rhinophores and the posterior part has six contractile gills inserted independently in the body. The pair of oral tentacles are constituted by a fine flexible membrane provided with large digital lobes. In a normal situation when the animal is crawling, the edges of its mantle are curled upwards creating a peripheral blister. If the animal is disturbed, it unfolds its edges and can swim through contractions and undulations of the body to move away from the disturbing element. Its common name, Spanish dancer, comes from this particular defense. Distribution and habitat This species is reported from the Red Sea, the Western Indian Ocean, French Polynesia and the Western Pacific, with different colour morphs in each region which are not differentiated by morphology or DNA barcodes. It likes rocky and coral reefs with many sponges and shelters from 1 to 50 meters deep.
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Biology During daytime, the Spanish dancer hides away from the light in the crevices of its natural habitat to only come out late at night. It feeds on various species of sponge. Like all nudibranchs, it is hermaphrodite and its bright red to pink egg ribbon has a spiral shape related to the size of the animal so relatively large. The Spanish Dancer consumes sponges from the family Halichondriidae. Once consumed, the Spanish Dancer derives a potent chemical that it can use as defense. Hexabranchus sanguineus then passes the defensive compounds obtained into its egg ribbons via macrolides, giving the physically defenseless egg ribbons a toxin defense. The latter is coveted by some other species of nudibranch as Favorinus tsuruganus or Favorinus japonicus. The emperor shrimp, Periclimenes imperator, is a commensal shrimp that is commonly found living on Hexabranchus sanguineus.
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An antimicrobial is an agent that kills microorganisms (microbicide) or stops their growth (bacteriostatic agent). Antimicrobial medicines can be grouped according to the microorganisms they act primarily against. For example, antibiotics are used against bacteria, and antifungals are used against fungi. They can also be classified according to their function. Antimicrobial medicines to treat infection are known as ⠀⠀antimicrobial chemotherapy, while antimicrobial drugs are used to prevent infection, which known as antimicrobial prophylaxis. The main classes of antimicrobial agents are disinfectants (non-selective agents, such as bleach), which kill a wide range of microbes on non-living surfaces to prevent the spread of illness, antiseptics (which are applied to living tissue and help reduce infection during surgery), and antibiotics (which destroy microorganisms within the body). The term antibiotic originally described only those formulations derived from living microorganisms but is now also applied to synthetic agents, such as sulfonamide's or fluoroquinolone's. Though the term used to be restricted to antibacterial, and is often used as a synonym for them by medical professionals and in medical literature, its context has broadened to include all antimicrobials. Antibacterial agents can be further subdivided into bactericidal agents, which kill bacteria, and bacteriostatic agents, which slow down or stall bacterial growth. In response, further advancements in antimicrobial technologies have resulted in solutions that can go beyond simply inhibiting microbial growth. Instead, certain types of porous media have been developed to kill microbes on contact. The misuse and overuse of antimicrobials in humans, animals and plants are the main drivers in the development of drug-resistant pathogens. It is estimated that bacterial antimicrobial resistance (AMR) was directly responsible for 1.27 million global deaths in 2019 and contributed to 4.95 million deaths. History Antimicrobial use has been common practice for at least 2000 years. Ancient Egyptians and ancient Greeks used specific molds and plant extracts to treat infection.
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In the 19th century, microbiologists such as Louis Pasteur and Jules Francois Joubert observed antagonism between some bacteria and discussed the merits of controlling these interactions in medicine. Louis Pasteur's work in fermentation and spontaneous generation led to the distinction between anaerobic and aerobic bacteria. The information garnered by Pasteur led Joseph Lister to incorporate antiseptic methods, such as sterilizing surgical tools and debriding wounds into surgical procedures. The implementation of these antiseptic techniques drastically reduced the number of infections and subsequent deaths associated with surgical procedures. Louis Pasteur's work in microbiology also led to the development of many vaccines for life-threatening diseases such as anthrax and rabies. On September 3, 1928, Alexander Fleming returned from a vacation and discovered that a Petri dish filled with Staphylococcus was separated into colonies due to the antimicrobial fungus Penicillium rubens. Fleming and his associates struggled to isolate the antimicrobial but referenced its therapeutic potential in 1929 in the British Journal of Experimental Pathology. In 1942, Howard Florey, Ernst Chain, and Edward Abraham used Fleming's work to purify and extract penicillin for medicinal uses earning them the 1945 Nobel Prize in Medicine. Chemical Antibacterials Antibacterials are used to treat bacterial infections. Antibiotics are classified generally as beta-lactams, macrolides, quinolones, tetracyclines or aminoglycosides. Their classification within these categories depends on their antimicrobial spectra, pharmacodynamics and chemical composition. Prolonged use of certain antibacterials can decrease the number of enteric bacteria, which may have a negative impact on health. Consumption of probiotics and healthy eating may help to replace destroyed gut flora. Stool transplants may be considered however for patients who are having difficulty recovering from prolonged antibiotic treatment, such as recurrent Clostridioides difficile infections.
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The discovery, development and use of antibacterials during the 20th century have reduced mortality from bacterial infections. The antibiotic era began with the therapeutic application of sulfonamide drugs in 1936, followed by a "golden" period of discovery from about 1945 to 1970, when a number of structurally diverse and highly effective agents were discovered and developed. Since 1980, the introduction of new antimicrobial agents for clinical use has declined, in part because of the enormous expense of developing and testing new drugs. In parallel, there has been an alarming increase in antimicrobial resistance of bacteria, fungi, parasites and some viruses to multiple existing agents. Antibacterials are among the most commonly used and misused drugs by physicians, for example, in viral respiratory tract infections. As a consequence of widespread and injudicious use of antibacterials, there has been an accelerated emergence of antibiotic-resistant pathogens, resulting in a serious threat to global public health. The resistance problem demands that a renewed effort be made to seek antibacterial agents effective against pathogenic bacteria resistant to current antibacterials. Possible strategies towards this objective include increased sampling from diverse environments and application of metagenomics to identify bioactive compounds produced by currently unknown and uncultured microorganisms as well as the development of small-molecule libraries customized for bacterial targets. Antifungals Antifungals are used to kill or prevent further growth of fungi. In medicine, they are used as a treatment for infections such as athlete's foot, ringworm and thrush and work by exploiting differences between mammalian and fungal cells. Unlike bacteria, both fungi and humans are eukaryotes. Thus, fungal and human cells are similar at the molecular level, making it more difficult to find a target for an antifungal drug to attack that does not also exist in the host organism. Consequently, there are often side effects to some of these drugs. Some of these side effects can be life-threatening if the drug is not used properly.
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As well as their use in medicine, antifungals are frequently sought after to control indoor mold in damp or wet home materials. Sodium bicarbonate (baking soda) blasted on to surfaces acts as an antifungal. Another antifungal solution applied after or without blasting by soda is a mix of hydrogen peroxide and a thin surface coating that neutralizes mold and encapsulates the surface to prevent spore release. Some paints are also manufactured with an added antifungal agent for use in high humidity areas such as bathrooms or kitchens. Other antifungal surface treatments typically contain variants of metals known to suppress mold growth e.g. pigments or solutions containing copper, silver or zinc. These solutions are not usually available to the general public because of their toxicity. Antivirals Antiviral drugs are a class of medication used specifically for treating viral infections. Like antibiotics, specific antivirals are used for specific viruses. They should be distinguished from viricides, which actively deactivate virus particles outside the body. Many antiviral drugs are designed to treat infections by retroviruses, including HIV. Important antiretroviral drugs include the class of protease inhibitors. Herpes viruses, best known for causing cold sores and genital herpes, are usually treated with the nucleoside analogue acyclovir. Viral hepatitis is caused by five unrelated hepatotropic viruses (A-E) and may be treated with antiviral drugs depending on the type of infection. Some influenza A and B viruses have become resistant to neuraminidase inhibitors such as oseltamivir, and the search for new substances continues. Antiparasitics Antiparasitics are a class of medications indicated for the treatment of infectious diseases such as leishmaniasis, malaria and Chagas disease, which are caused by parasites such as nematodes, cestodes, trematodes and infectious protozoa. Antiparasitic medications include metronidazole, iodoquinol and albendazole. Like all therapeutic antimicrobials, they must kill the infecting organism without serious damage to the host. Broad-spectrum therapeutics Broad-spectrum therapeutics are active against multiple classes of pathogens. Such therapeutics have been suggested as potential emergency treatments for pandemics. Non-pharmaceutical
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A wide range of chemical and natural compounds are used as antimicrobials. Organic acids and their salts are used widely in food products, e.g. lactic acid, citric acid, acetic acid, either as ingredients or as disinfectants. For example, beef carcasses often are sprayed with acids, and then rinsed or steamed, to reduce the prevalence of Escherichia coli. Heavy metal cations such as Hg2+ and Pb2+ have antimicrobial activities, but can be toxic. In recent years, the antimicrobial activity of coordination compounds has been investigated. Traditional herbalists used plants to treat infectious disease. Many of these plants have been investigated scientifically for antimicrobial activity, and some plant products have been shown to inhibit the growth of pathogenic microorganisms. A number of these agents appear to have structures and modes of action that are distinct from those of the antibiotics in current use, suggesting that cross-resistance with agents already in use may be minimal. Copper Copper-alloy surfaces have natural intrinsic antimicrobial properties and can kill microorganisms such as E. coli and Staphylococcus. The United States Environmental Protection Agency approved the registration of antimicrobial copper alloy surfaces for use in addition to regular cleaning and disinfection to control infections. Antimicrobial copper alloys are being installed in some healthcare facilities and subway transit systems as a public hygienic measure. Copper nanoparticles are attracting interest for the intrinsic antimicrobial behaviours. Essential oils Many essential oils included in herbal pharmacopoeias are claimed to possess antimicrobial activity, with the oils of bay, cinnamon, clove and thyme reported to be the most potent in studies with foodborne bacterial pathogens. Coconut oil is also known for its antimicrobial properties. Active constituents include terpenoids and secondary metabolites. Despite their prevalent use in alternative medicine, essential oils have seen limited use in mainstream medicine. While 25 to 50% of pharmaceutical compounds are plant-derived, none are used as antimicrobials, though there has been increased research in this direction. Barriers to increased usage in mainstream medicine include poor regulatory oversight and quality control, mislabeled or misidentified products, and limited modes of delivery.
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Antimicrobial pesticides According to the U.S. Environmental Protection Agency (EPA), and defined by the Federal Insecticide, Fungicide, and Rodenticide Act, antimicrobial pesticides are used to control growth of microbes through disinfection, sanitation, or reduction of development and to protect inanimate objects, industrial processes or systems, surfaces, water, or other chemical substances from contamination, fouling, or deterioration caused by bacteria, viruses, fungi, protozoa, algae, or slime. The EPA monitors products, such as disinfectants/sanitizers for use in hospitals or homes, to ascertain efficacy. Products that are meant for public health are therefore under this monitoring system, including products used for drinking water, swimming pools, food sanitation, and other environmental surfaces. These pesticide products are registered under the premise that, when used properly, they do not demonstrate unreasonable side effects to humans or the environment. Even once certain products are on the market, the EPA continues to monitor and evaluate them to make sure they maintain efficacy in protecting public health. Public health products regulated by the EPA are divided into three categories: Disinfectants: Destroy or inactivate microorganisms (bacteria, fungi, viruses,) but may not act as sporicides (as those are the most difficult form to destroy). According to efficacy data, the EPA will classify a disinfectant as limited, general/ broad spectrum, or as a hospital disinfectant. Sanitizers: Reduce the number of microorganisms, but may not kill or eliminate all of them. Sterilizers (Sporicides): Eliminate all bacteria, fungi, spores, and viruses. Antimicrobial pesticide safety Antimicrobial pesticides have the potential to be a major factor in drug resistance. Organizations such as the World Health Organization call for significant reduction in their use globally to combat this. According to a 2010 Centers for Disease Control and Prevention report, health-care workers can take steps to improve their safety measures against antimicrobial pesticide exposure. Workers are advised to minimize exposure to these agents by wearing personal protective equipment such as gloves and safety glasses. Additionally, it is important to follow the handling instructions properly, as that is how the EPA has deemed them as safe to use. Employees should be educated about the health hazards and encouraged to seek medical care if exposure occurs. Ozone
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Ozone can kill microorganisms in air, water and process equipment and has been used in settings such as kitchen exhaust ventilation, garbage rooms, grease traps, biogas plants, wastewater treatment plants, textile production, breweries, dairies, food and hygiene production, pharmaceutical industries, bottling plants, zoos, municipal drinking-water systems, swimming pools and spas, and in the laundering of clothes and treatment of in–house mold and odors. Antimicrobial scrubs Antimicrobial scrubs can reduce the accumulation of odors and stains on scrubs, which in turn improves their longevity. These scrubs also come in a variety of colors and styles. As antimicrobial technology develops at a rapid pace, these scrubs are readily available, with more advanced versions hitting the market every year. These bacteria could then be spread to office desks, break rooms, computers, and other shared technology. This can lead to outbreaks and infections like methicillin-resistant staphylococcus aureus, treatments for which cost the healthcare industry $20 billion a year.
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Halogens Elements such as chlorine, iodine, fluorine, and bromine are nonmetallic in nature and constitute the halogen family. Each of these halogens have a different antimicrobial effect that is influenced by various factors such as pH, temperature, contact time, and type of microorganism. Chlorine and iodine are the two most commonly used antimicrobials. Chlorine is extensively used as a disinfectant in the water treatment plants, drug, and food industries. In wastewater treatment plants, chlorine is widely used as a disinfectant. It oxidizes soluble contaminants and kills bacteria and viruses. It is also highly effective against bacterial spores. The mode of action is by breaking the bonds present in these microorganisms. When a bacterial enzyme comes in contact with a compound containing chlorine, the hydrogen atom in that molecule gets displaced and is replaced with chlorine. This in turn changes the enzyme function which ultimately leads to the death of the bacterium. Iodine is most commonly used for sterilization and wound cleaning. The three major antimicrobial compounds containing iodine are alcohol-iodine solution, an aqueous solution of iodine, and iodophors. Iodophors are more bactericidal and are used as antiseptics as they are less irritating when applied to the skin. Bacterial spores on the other hand cannot be killed by iodine, but they can be inhibited by iodophors. The growth of microorganisms is inhibited when iodine penetrates into the cells and oxidizes proteins, genetic material, and fatty acids. Bromine is also an effective antimicrobial that is used in water treatment plants. When mixed with chlorine it is highly effective against bacterial spores such as S. faecalis.
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Alcohols Alcohols are commonly used as disinfectants and antiseptics. Alcohols kill vegetative bacteria, most viruses and fungi. Ethyl alcohol, n-propanol and isopropyl alcohol are the most commonly used antimicrobial agents. Methanol is also a disinfecting agent but is not generally used as it is highly poisonous. Escherichia coli, Salmonella, and Staphylococcus aureus are a few bacteria whose growth can be inhibited by alcohols. Alcohols have a high efficiency against enveloped viruses (60–70% ethyl alcohol) 70% isopropyl alcohol or ethanol are highly effective as an antimicrobial agent. In the presence of water, 70% alcohol causes coagulation of the proteins thus inhibiting microbial growth. Alcohols are not quite efficient when it comes to spores. The mode of action is by denaturing the proteins. Alcohols interfere with the hydrogen bonds present in the protein structure. Alcohols also dissolve the lipid membranes that are present in microorganisms. Disruption of the cell membrane is another property of alcohols that aids in cell death. Alcohols are cheap and effective antimicrobials. They are widely used in the pharmaceutical industry.  Alcohols are commonly used in hand sanitizers, antiseptics, and disinfectants.
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Phenol and Phenolic compounds Phenol, also known as carbolic acid, was one of the first chemicals which was used as an antimicrobial agent. It has high antiseptic properties. It is bacteriostatic at concentrations of 0.1%–1% and is bactericidal/fungicidal at 1%–2%. A 5% solution kills anthrax spores in 48 hr. Phenols are most commonly used in oral mouth washes and household cleaning agents. They are active against a wide range of bacteria, fungi and viruses.  Today phenol derivatives such as thymol and cresol are used because they are less toxic compared to phenol. These phenolic compounds have a benzene ring along with the –OH group incorporated into their structures. They have a higher antimicrobial activity. These compounds inhibit microbial growth by precipitating proteins which lead to their denaturation and by penetrating into the cell membrane of microorganisms and disrupting it. Phenolic compounds can also deactivate enzymes and damage the amino acids in microbial cells. Phenolics such as fentichlore, an antibacterial and antifungal agent, are used as an oral treatment for fungal infections. Trischlosan is highly effective against both gram-positive and gram-negative bacteria. Hexachlorophene (Bisphenol) is used as a surfactant. It is widely used in soaps, handwashes, and skin products because of its antiseptic properties. It is also used as a sterilizing agent. Cresol is an effective antimicrobial and is widely used in mouthwashes and cough drops. Phenolics have high antimicrobial activity against bacteria such as Staphylococcus epidermidis and Pseudomonas aeruginosa. 2-Phenylphenol-water solutions are used in immersion treatments of fruit for packing. (It is not used on the packing materials however.) Ihloff and Kalitzki 1961 find a small but measurable amount remains in the skin of fruits processed in this manner.
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Aldehydes Aldehydes are highly effective against bacteria, fungi, and viruses.  Aldehydes inhibit bacterial growth by disrupting the outer membrane. They are used in the disinfection and sterilization of surgical instruments. As they are highly toxic, they are not used in antiseptics. Currently, only three aldehyde compounds are of widespread practical use as disinfectant biocides, namely glutaraldehyde, formaldehyde, and ortho-phthalaldehyde (OPA) despite the demonstration that many other aldehydes possess good antimicrobial activity. However, due to its long contact time other disinfectants are commonly preferred. Physical Heat Microorganisms have a minimum temperature, an optimum, and a maximum temperature for growth. High temperature as well as low temperatures are used as physical agents of control. Different organisms show different degrees of resistance or susceptibility to heat or temperature, some organisms such as bacterial endospore are more resistant while vegetative cells are less resistant and are easily killed at lower temperatures. Another method that involves the use of heat to kill microorganisms is fractional sterilization. This process involves the exposure to a temperature of 100 degrees Celsius for an hour per day for several days. Fractional sterilization is also called tyndallization. Bacterial endospores can be killed using this method. Both dry and moist heat are effective in eliminating microbial life. For example, jars used to store preserves such as jam can be sterilized by heating them in a conventional oven. Heat is also used in pasteurization, a method for slowing the spoilage of foods such as milk, cheese, juices, wines and vinegar. Such products are heated to a certain temperature for a set period of time, which greatly reduces the number of harmful microorganisms. Low temperature is also used to inhibit microbial activity by slowing down microbial metabolism. Radiation Foods are often irradiated to kill harmful pathogens. There are two types of radiations that are used to inhibit the growth of microorganisms – ionizing and non-ionizing radiations. Common sources of radiation used in food sterilization include cobalt-60 (a gamma emitter), electron beams and . Ultraviolet light is also used to disinfect drinking water, both in small-scale personal-use systems and larger-scale community water purification systems.
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Desiccation Desiccation is also known as dehydration.  It is the state of extreme dryness or the process of extreme drying. Some microorganisms like bacteria, yeasts and molds require water for their growth. Desiccation dries up the water content thus inhibiting microbial growth. On the availability of water, the bacteria resume their growth, thus desiccation does not completely inhibit bacterial growth. The instrument used to carry out this process is called a desiccator. This process is widely used in the food industry and is an efficient method for food preservation. Desiccation is also largely used in the pharmaceutical industry to store vaccines and other products. Antimicrobial surfaces Antimicrobial surfaces are designed to either inhibit the ability of microorganisms to grow or damaging them by chemical (copper toxicity) or physical processes (micro/nano-pillars to rupture cell walls). These surfaces are especially important for the healthcare industry. Designing effective antimicrobial surfaces demands an in-depth understanding of the initial microbe-surface adhesion mechanisms. Molecular dynamics simulation and time-lapse imaging are typically used to investigate these mechanisms. Osmotic pressure Osmotic pressure is the pressure required to prevent a solvent from passing from a region of high concentration to a region of low concentration through a semipermeable membrane.  When the concentration of dissolved materials or solute is higher inside the cell than it is outside, the cell is said to be in a hypotonic environment and water will flow into the cell.When the bacteria is placed in hypertonic solution, it causes plasmolysis or cell shrinking, similarly in hypotonic solution, bacteria undergoes plasmotysis or turgid state. This plasmolysis and plasmotysis kills bacteria because it causes change in osmotic pressure. Antimicrobial resistance Antimicrobial resistance The misuse and overuse of antimicrobials in humans, animals and plants are the main drivers in the development of drug-resistant pathogens. It is estimated that bacterial antimicrobial resistance (AMR) was directly responsible for 1.27 million global deaths in 2019 and contributed to 4.95 million deaths.
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The roller is an agricultural tool used for flattening land or breaking up large clumps of soil, especially after ploughing or disc harrowing. Typically, rollers are pulled by tractors or, prior to mechanisation, a team of animals such as horses or oxen. As well as for agricultural purposes, rollers are used on cricket pitches and residential lawn areas. Flatter land makes subsequent weed control and harvesting easier, and rolling can help to reduce moisture loss from cultivated soil. On lawns, rolling levels the land for mowing and compacts the soil surface. Rollers may be weighted in different ways. For many uses a heavy roller is used. These may consist of one or more cylinders made of thick steel, a thinner steel cylinder filled with concrete, or a cylinder filled with water. A water-filled roller has the advantage that the water may be drained out for lighter use or for transport. In frost-prone areas a water filled roller must be drained for winter storage to avoid breakage due to the expansion for water as it turns to ice. Designs One-piece versus segmented On tilled soil a one-piece roller has the disadvantage that when turning corners the outer end of the roller has to rotate much faster than the inner end, forcing one or both ends to skid. A one-piece roller turned on soft ground will skid up a heap of soil at the outer radius, leaving heaps, which is counter-productive. Rollers are often made in two or three sections to reduce this problem, and the Cambridge roller overcomes it altogether by mounting many small segments onto one axle so that they can each rotate at local ground-speed. Smooth versus ridged The surface of rollers may be smooth, or it may be textured to help break up soil or to groove the final surface to reduce scouring from rain. Each segment of a Cambridge roller has a rib around its edge for this purpose. The name cultipacker is often used for such ridged types, especially in the United States. Uses
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Farming use Rollers are a secondary tillage tool used for flattening land or breaking up large clumps of soil, especially after ploughing or disc harrowing. Rollers are typically pulled by tractors today. Before mechanised agriculture, a team of working animals such as horses or oxen provided the power. Animal power is still used today in some contexts, such as on Amish farms in the United States and in regions of Asia where draft oxen are still widely used. Rollers prepare optimal seedbeds by making them as flat as is practical and moderately firmed. Flatness is important at planting because it is the only practical way to control average seed planting depth without laborious hand planting of each seed; it is not practical to follow an instruction of (for example) 1-cm planting depth if the contour of the seedbed varies by 2 cm or more between adjacent spots. This is why breaking up of even small clods/lumps, and well-leveled spreading of soil, is important at planting time. Flatter land also makes subsequent weed control and harvesting easier. For example, in mechanical weed control, controlling cultivator tooth depth is practical only with a decently flat soil contour, and in combining, controlling combine head height is practical only with a decently flat soil contour. Rolling is also believed to help reduce moisture loss from cultivated soil. Ganging and trailing Rollers may be ganged to increase the width of each pass/swath. Rollers may be trailed after other equipment such as ploughs, disc harrows, or mowers. Cricket pitch use In cricket, rollers are used to make the pitch flat and less dangerous for batsmen. Several size rollers have been used in the history of cricket, from light rollers that were used in the days of uncovered pitches and at some stages during the 1950s to make batting less easy, to the modern “heavy roller” universally used in top-class cricket today. Regulations permit a pitch only to be rolled at the commencement of each innings or day’s play, but this has still had a massive influence on the game by eliminating the shooters that were ubiquitous on all but light soils before heavy rollers were used. Heavy rollers have sometimes been criticised for making batting too easy and for reducing the rate at which pitches dry out after rain in the cool English climate. Lawn use
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Lawn rollers are designed to even out or firm up the lawn surface, especially in climates where heaving causes the lawn to be lumpy. Heaving may result when the ground freezes and thaws many times over winter. Where this occurs, gardeners are advised to give the lawn a light rolling with a lawn roller in the spring. Clay or wet soils should not be rolled as they become compacted.
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The inverse second or reciprocal second (s−1), also called per second, is a unit defined as the multiplicative inverse of the second (a unit of time). It is applicable for physical quantities of dimension reciprocal time, such as frequency and strain rate. It is dimensionally equivalent to: hertz (Hz), historically known as cycles per second – the SI unit for frequency and rotational frequency becquerel (Bq) – the SI unit for the rate of occurrence of aperiodic or stochastic radionuclide events baud (Bd) – the unit for symbol rate over a communication link bit per second (bit/s) – the unit of bit rate However, the special names and symbols above for s−1 are recommend for clarity. Reciprocal second should not be confused with radian per second (rad⋅s−1), the SI unit for angular frequency and angular velocity. As the radian is a dimensionless unit, radian per second is dimensionally consistent with reciprocal second. However, they are used for different kinds of quantity, frequency and angular frequency, whose numerical value differs by 2. The inverse minute or reciprocal minute (min−1), also called per minute, is 60−1 s−1, as 1 min = 60 s; it is used in quantities of type "counts per minute", such as: Actions per minute Beats per minute Counts per minute Revolutions per minute (rpm) Words per minute Inverse square second (s−2) is involved in the units of linear acceleration, angular acceleration, and rotational acceleration.
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Sugar alcohols (also called polyhydric alcohols, polyalcohols, alditols or glycitols) are organic compounds, typically derived from sugars, containing one hydroxyl group attached to each carbon atom. They are white, water-soluble solids that can occur naturally or be produced industrially by hydrogenating sugars. Since they contain multiple groups, they are classified as polyols. Sugar alcohols are used widely in the food industry as thickeners and sweeteners. In commercial foodstuffs, sugar alcohols are commonly used in place of table sugar (sucrose), often in combination with high-intensity artificial sweeteners, in order to offset their low sweetness. Xylitol and sorbitol are popular sugar alcohols in commercial foods. Structure Sugar alcohols have the general formula . In contrast, sugars have two fewer hydrogen atoms, for example, or . Like their parent sugars, sugar alcohols exist in diverse chain length. Most have five- or six-carbon chains, because they are derived respectively from pentoses (five-carbon sugars) and hexoses (six-carbon sugars), which are the more common sugars. They have one −OH group attached to each carbon. They are further differentiated by the relative orientation (stereochemistry) of these −OH groups. Unlike sugars, which tend to exist as rings, sugar alcohols do not, although they can be dehydrated to give cyclic ethers (e.g. sorbitan can be dehydrated to isosorbide). Production Sugar alcohols can be, and often are, produced from renewable resources. Particular feedstocks are starch, cellulose and hemicellulose; the main conversion technologies use as the reagent: hydrogenolysis, i.e. the cleavage of single bonds, converting polymers to smaller molecules, and hydrogenation of double bonds, converting sugars to sugar alcohols. Sorbitol and mannitol Mannitol is no longer obtained from natural sources; currently, sorbitol and mannitol are obtained by hydrogenation of sugars, using Raney nickel catalysts. The conversion of glucose and mannose to sorbitol and mannitol is given as Erythritol Erythritol is obtained by the fermentation of glucose and sucrose.
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Health effects Sugar alcohols do not contribute to tooth decay; in fact, xylitol deters tooth decay. Sugar alcohols are absorbed at 50% of the rate of sugars, resulting in less of an effect on blood sugar levels as measured by comparing their effect to sucrose using the glycemic index. Common sugar alcohols Ethylene glycol (2-carbon) Glycerol (3-carbon) Erythritol (4-carbon) Threitol (4-carbon) Arabitol (5-carbon) Xylitol (5-carbon) Ribitol (5-carbon) Mannitol (6-carbon) Sorbitol (6-carbon) Galactitol (6-carbon) Fucitol (6-carbon) Iditol (6-carbon) Inositol (6-carbon; a cyclic sugar alcohol) Volemitol (7-carbon) Isomalt (12-carbon) Maltitol (12-carbon) Lactitol (12-carbon) Maltotriitol (18-carbon) Maltotetraitol (24-carbon) Polyglycitol Both disaccharides and monosaccharides can form sugar alcohols; however, sugar alcohols derived from disaccharides (e.g. maltitol and lactitol) are not entirely hydrogenated because only one aldehyde group is available for reduction. Sugar alcohols as food additives This table presents the relative sweetness and food energy of the most widely used sugar alcohols. Despite the variance in food energy content of sugar alcohols, the European Union's labeling requirements assign a blanket value of 2.4 kcal/g to all sugar alcohols. Characteristics As a group, sugar alcohols are not as sweet as sucrose, and they have slightly less food energy than sucrose. Their flavor is similar to sucrose, and they can be used to mask the unpleasant aftertastes of some high-intensity sweeteners. Sugar alcohols are not metabolized by oral bacteria, and so they do not contribute to tooth decay. They do not brown or caramelize when heated.
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In addition to their sweetness, some sugar alcohols can produce a noticeable cooling sensation in the mouth when highly concentrated, for instance in sugar-free hard candy or chewing gum. This happens, for example, with the crystalline phase of sorbitol, erythritol, xylitol, mannitol, lactitol and maltitol. The cooling sensation is due to the dissolution of the sugar alcohol being an endothermic (heat-absorbing) reaction, one with a strong heat of solution. Absorption from the small intestine Sugar alcohols are usually incompletely absorbed into the blood stream from the small intestine which generally results in a smaller change in blood glucose than "regular" sugar (sucrose). This property makes them popular sweeteners among diabetics and people on low-carbohydrate diets. As an exception, erythritol is actually absorbed in the small intestine and excreted unchanged through urine, so it contributes no calories even though it is rather sweet. Side effects Like many other incompletely digestible substances, overconsumption of sugar alcohols can lead to bloating, diarrhea and flatulence because they are not fully absorbed in the small intestine. Some individuals experience such symptoms even in a single-serving quantity. With continued use, most people develop a degree of tolerance to sugar alcohols and no longer experience these symptoms.
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In graph theory, a vertex cover (sometimes node cover) of a graph is a set of vertices that includes at least one endpoint of every edge of the graph. In computer science, the problem of finding a minimum vertex cover is a classical optimization problem. It is NP-hard, so it cannot be solved by a polynomial-time algorithm if P ≠ NP. Moreover, it is hard to approximate – it cannot be approximated up to a factor smaller than 2 if the unique games conjecture is true. On the other hand, it has several simple 2-factor approximations. It is a typical example of an NP-hard optimization problem that has an approximation algorithm. Its decision version, the vertex cover problem, was one of Karp's 21 NP-complete problems and is therefore a classical NP-complete problem in computational complexity theory. Furthermore, the vertex cover problem is fixed-parameter tractable and a central problem in parameterized complexity theory. The minimum vertex cover problem can be formulated as a half-integral, linear program whose dual linear program is the maximum matching problem. Vertex cover problems have been generalized to hypergraphs, see Vertex cover in hypergraphs. Definition Formally, a vertex cover of an undirected graph is a subset of such that , that is to say it is a set of vertices where every edge has at least one endpoint in the vertex cover . Such a set is said to cover the edges of . The upper figure shows two examples of vertex covers, with some vertex cover marked in red. A minimum vertex cover is a vertex cover of smallest possible size. The vertex cover number is the size of a minimum vertex cover, i.e. . The lower figure shows examples of minimum vertex covers in the previous graphs. Examples The set of all vertices is a vertex cover. The endpoints of any maximal matching form a vertex cover. The complete bipartite graph has a minimum vertex cover of size . Properties A set of vertices is a vertex cover if and only if its complement is an independent set. Consequently, the number of vertices of a graph is equal to its minimum vertex cover number plus the size of a maximum independent set. Computational problem The minimum vertex cover problem is the optimization problem of finding a smallest vertex cover in a given graph. INSTANCE: Graph OUTPUT: Smallest number such that has a vertex cover of size . If the problem is stated as a decision problem, it is called the vertex cover problem: INSTANCE: Graph and positive integer . QUESTION: Does have a vertex cover of size at most ?
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The vertex cover problem is an NP-complete problem: it was one of Karp's 21 NP-complete problems. It is often used in computational complexity theory as a starting point for NP-hardness proofs. ILP formulation Assume that every vertex has an associated cost of . The (weighted) minimum vertex cover problem can be formulated as the following integer linear program (ILP). {| | minimize | colspan="2" | |    | (minimize the total cost) |- | subject to | | for all | | (cover every edge of the graph), |- | | | for all . | | (every vertex is either in the vertex cover or not) |} This ILP belongs to the more general class of ILPs for covering problems. The integrality gap of this ILP is , so its relaxation (allowing each variable to be in the interval from 0 to 1, rather than requiring the variables to be only 0 or 1) gives a factor- approximation algorithm for the minimum vertex cover problem. Furthermore, the linear programming relaxation of that ILP is half-integral, that is, there exists an optimal solution for which each entry is either 0, 1/2, or 1. A 2-approximate vertex cover can be obtained from this fractional solution by selecting the subset of vertices whose variables are nonzero. Exact evaluation The decision variant of the vertex cover problem is NP-complete, which means it is unlikely that there is an efficient algorithm to solve it exactly for arbitrary graphs. NP-completeness can be proven by reduction from 3-satisfiability or, as Karp did, by reduction from the clique problem. Vertex cover remains NP-complete even in cubic graphs and even in planar graphs of degree at most 3. For bipartite graphs, the equivalence between vertex cover and maximum matching described by Kőnig's theorem allows the bipartite vertex cover problem to be solved in polynomial time. For tree graphs, an algorithm finds a minimal vertex cover in polynomial time by finding the first leaf in the tree and adding its parent to the minimal vertex cover, then deleting the leaf and parent and all associated edges and continuing repeatedly until no edges remain in the tree. Fixed-parameter tractability
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An exhaustive search algorithm can solve the problem in time 2knO(1), where k is the size of the vertex cover. Vertex cover is therefore fixed-parameter tractable, and if we are only interested in small k, we can solve the problem in polynomial time. One algorithmic technique that works here is called bounded search tree algorithm, and its idea is to repeatedly choose some vertex and recursively branch, with two cases at each step: place either the current vertex or all its neighbours into the vertex cover. The algorithm for solving vertex cover that achieves the best asymptotic dependence on the parameter runs in time . The klam value of this time bound (an estimate for the largest parameter value that could be solved in a reasonable amount of time) is approximately 190. That is, unless additional algorithmic improvements can be found, this algorithm is suitable only for instances whose vertex cover number is 190 or less. Under reasonable complexity-theoretic assumptions, namely the exponential time hypothesis, the running time cannot be improved to 2o(k), even when is . However, for planar graphs, and more generally, for graphs excluding some fixed graph as a minor, a vertex cover of size k can be found in time , i.e., the problem is subexponential fixed-parameter tractable. This algorithm is again optimal, in the sense that, under the exponential time hypothesis, no algorithm can solve vertex cover on planar graphs in time . Approximate evaluation One can find a factor-2 approximation by repeatedly taking both endpoints of an edge into the vertex cover, then removing them from the graph. Put otherwise, we find a maximal matching M with a greedy algorithm and construct a vertex cover C that consists of all endpoints of the edges in M. In the following figure, a maximal matching M is marked with red, and the vertex cover C is marked with blue.
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The set C constructed this way is a vertex cover: suppose that an edge e is not covered by C; then M ∪ {e} is a matching and e ∉ M, which is a contradiction with the assumption that M is maximal. Furthermore, if e = {u, v} ∈ M, then any vertex cover – including an optimal vertex cover – must contain u or v (or both); otherwise the edge e is not covered. That is, an optimal cover contains at least one endpoint of each edge in M; in total, the set C is at most 2 times as large as the optimal vertex cover. This simple algorithm was discovered independently by Fanica Gavril and Mihalis Yannakakis. More involved techniques show that there are approximation algorithms with a slightly better approximation factor. For example, an approximation algorithm with an approximation factor of is known. The problem can be approximated with an approximation factor in - dense graphs. Inapproximability No better constant-factor approximation algorithm than the above one is known. The minimum vertex cover problem is APX-complete, that is, it cannot be approximated arbitrarily well unless P = NP. Using techniques from the PCP theorem, Dinur and Safra proved in 2005 that minimum vertex cover cannot be approximated within a factor of 1.3606 for any sufficiently large vertex degree unless P = NP. Later, the factor was improved to for any . Moreover, if the unique games conjecture is true then minimum vertex cover cannot be approximated within any constant factor better than 2. Although finding the minimum-size vertex cover is equivalent to finding the maximum-size independent set, as described above, the two problems are not equivalent in an approximation-preserving way: The Independent Set problem has no constant-factor approximation unless P = NP. Pseudocode APPROXIMATION-VERTEX-COVER(G) C = ∅ E'= G.E while E' ≠ ∅: let (u, v) be an arbitrary edge of E' C = C ∪ {u, v} remove from E' every edge incident on either u or v return C
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Applications Vertex cover optimization serves as a model for many real-world and theoretical problems. For example, a commercial establishment interested in installing the fewest possible closed circuit cameras covering all hallways (edges) connecting all rooms (nodes) on a floor might model the objective as a vertex cover minimization problem. The problem has also been used to model the elimination of repetitive DNA sequences for synthetic biology and metabolic engineering applications.
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In mathematics and logic, a higher-order logic (abbreviated HOL) is a form of logic that is distinguished from first-order logic by additional quantifiers and, sometimes, stronger semantics. Higher-order logics with their standard semantics are more expressive, but their model-theoretic properties are less well-behaved than those of first-order logic. The term "higher-order logic" is commonly used to mean higher-order simple predicate logic. Here "simple" indicates that the underlying type theory is the theory of simple types, also called the simple theory of types. Leon Chwistek and Frank P. Ramsey proposed this as a simplification of the complicated and clumsy ramified theory of types specified in the Principia Mathematica by Alfred North Whitehead and Bertrand Russell. Simple types is sometimes also meant to exclude polymorphic and dependent types. Quantification scope First-order logic quantifies only variables that range over individuals; second-order logic, also quantifies over sets; third-order logic also quantifies over sets of sets, and so on. Higher-order logic is the union of first-, second-, third-, ..., nth-order logic; i.e., higher-order logic admits quantification over sets that are nested arbitrarily deeply. Semantics There are two possible semantics for higher-order logic.
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In the standard or full semantics, quantifiers over higher-type objects range over all possible objects of that type. For example, a quantifier over sets of individuals ranges over the entire powerset of the set of individuals. Thus, in standard semantics, once the set of individuals is specified, this is enough to specify all the quantifiers. HOL with standard semantics is more expressive than first-order logic. For example, HOL admits categorical axiomatizations of the natural numbers, and of the real numbers, which are impossible with first-order logic. However, by a result of Kurt Gödel, HOL with standard semantics does not admit an effective, sound, and complete proof calculus. The model-theoretic properties of HOL with standard semantics are also more complex than those of first-order logic. For example, the Löwenheim number of second-order logic is already larger than the first measurable cardinal, if such a cardinal exists. The Löwenheim number of first-order logic, in contrast, is ℵ0, the smallest infinite cardinal. In Henkin semantics, a separate domain is included in each interpretation for each higher-order type. Thus, for example, quantifiers over sets of individuals may range over only a subset of the powerset of the set of individuals. HOL with these semantics is equivalent to many-sorted first-order logic, rather than being stronger than first-order logic. In particular, HOL with Henkin semantics has all the model-theoretic properties of first-order logic, and has a complete, sound, effective proof system inherited from first-order logic. Properties Higher-order logics include the offshoots of Church's simple theory of types and the various forms of intuitionistic type theory. Gérard Huet has shown that unifiability is undecidable in a type-theoretic flavor of third-order logic, that is, there can be no algorithm to decide whether an arbitrary equation between second-order (let alone arbitrary higher-order) terms has a solution. Up to a certain notion of isomorphism, the powerset operation is definable in second-order logic. Using this observation, Jaakko Hintikka established in 1955 that second-order logic can simulate higher-order logics in the sense that for every formula of a higher-order logic, one can find an equisatisfiable formula for it in second-order logic.
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The term "higher-order logic" is assumed in some context to refer to classical higher-order logic. However, modal higher-order logic has been studied as well. According to several logicians, Gödel's ontological proof is best studied (from a technical perspective) in such a context.
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Signs and symptoms are diagnostic indications of an illness, injury, or condition. Signs are objective and externally observable; symptoms are a person's reported subjective experiences. A sign for example may be a higher or lower temperature than normal, raised or lowered blood pressure or an abnormality showing on a medical scan. A symptom is something out of the ordinary that is experienced by an individual such as feeling feverish, a headache or other pains in the body, which occur as the body's immune system fights off an infection. Signs and symptoms Signs A medical sign is an objective observable indication of a disease, injury, or medical condition that may be detected during a physical examination. These signs may be visible, such as a rash or bruise, or otherwise detectable such as by using a stethoscope or taking blood pressure. Medical signs, along with symptoms, help in forming a diagnosis. Some examples of signs are nail clubbing of either the fingernails or toenails or an abnormal gait. Symptoms A symptom is something felt or experienced, such as pain or dizziness. Signs and symptoms are not mutually exclusive, for example a subjective feeling of fever can be noted as sign by using a thermometer that registers a high reading. The CDC lists various diseases by their signs and symptoms such as for measles which includes a high fever, conjunctivitis, and cough, followed a few days later by the measles rash. Cardinal signs and symptoms Cardinal signs and symptoms are very specific even to the point of being pathognomonic. A cardinal sign or cardinal symptom can also refer to the major sign or symptom of a disease. Abnormal reflexes can indicate problems with the nervous system. Signs and symptoms are also applied to physiological states outside the context of disease, as for example when referring to the signs and symptoms of pregnancy, or the symptoms of dehydration. Sometimes a disease may be present without showing any signs or symptoms when it is known as being asymptomatic. The disorder may be discovered through tests including scans. An infection may be asymptomatic but still be transmissible. Syndrome Signs and symptoms are often non-specific, but some combinations can be suggestive of certain diagnoses, helping to narrow down what may be wrong. A particular set of characteristic signs and symptoms that may be associated with a disorder is known as a syndrome. Related terms
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Symptomatic When a disease is evidenced by symptoms it is known as symptomatic. There are many conditions including subclinical infections that display no symptoms, and these are termed asymptomatic. Signs and symptoms may be mild or severe, brief or longer-lasting when they may become reduced (remission), or then recur (relapse or recrudescence) known as a flare-up. A flare-up may show more severe symptoms. The term chief complaint, also "presenting problem", is used to describe the initial concern of an individual when seeking medical help, and once this is clearly noted a history of the present illness may be taken. The symptom that ultimately leads to a diagnosis is called a cardinal symptom. Some symptoms can be misleading as a result of referred pain, where for example a pain in the right shoulder may be due to an inflamed gallbladder and not to presumed muscle strain. Prodrome Many diseases have an early prodromal stage where a few signs and symptoms may suggest the presence of a disorder before further specific symptoms may emerge. Measles for example has a prodromal presentation that includes a hacking cough, fever, and Koplik's spots in the mouth. Over half of migraine episodes have a prodromal phase. Schizophrenia has a notable prodromal stage, as has dementia. Nonspecific symptoms Some symptoms are specific, that is, they are associated with a single, specific medical condition. Nonspecific symptoms, sometimes also called equivocal symptoms, are not specific to a particular condition. They include unexplained weight loss, headache, pain, fatigue, loss of appetite, night sweats, and malaise. A group of three particular nonspecific symptoms – fever, night sweats, and weight loss – over a period of six months are termed B symptoms associated with lymphoma and indicate a poor prognosis. Other sub-types of symptoms include: constitutional or general symptoms, which affect general well-being or the whole body, such as a fever; concomitant symptoms, which are symptoms that occur at the same time as the primary symptom; prodromal symptoms, which are the first symptoms of an bigger set of problems; delayed symptoms, which happen some time after the trigger; and objective symptoms, which are symptoms whose existence can be observed and confirmed by a healthcare provider.
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Vital signs Vital signs are the four signs that can give an immediate measurement of the body's overall functioning and health status. They are temperature, heart rate, breathing rate, and blood pressure. The ranges of these measurements vary with age, weight, gender and with general health. A digital application has been developed for use in clinical settings that measures three of the vital signs (not temperature) using just a smartphone, and has been approved by NHS England. The application is registered as Lifelight First, and Lifelight Home is under development (2020) for monitoring-use by people at home using just the camera on their smartphone or tablet. This will additionally measure oxygen saturation and atrial fibrillation. Other devices are then not needed. Syndromes Many conditions are indicated by a group of known signs, or signs and symptoms. These can be a group of three known as a triad; a group of four ("tetrad"); or a group of five ("pentad"). An example of a triad is Meltzer's triad presenting purpura a rash, arthralgia painful joints, and myalgia painful and weak muscles. Meltzer's triad indicates the condition cryoglobulinemia. Huntington's disease is a neurodegenerative disease that is characterized by a triad of motor, cognitive, and psychiatric signs and symptoms. A large number of these groups that can be characteristic of a particular disease are known as a syndrome. Noonan syndrome for example, has a diagnostic set of unique facial and musculoskeletal features. Some syndromes such as nephrotic syndrome may have a number of underlying causes that are all related to diseases that affect the kidneys. Sometimes a child or young adult may have symptoms suggestive of a genetic disorder that cannot be identified even after genetic testing. In such cases the term SWAN (syndrome without a name) may be used. Often a diagnosis may be made at some future point when other more specific symptoms emerge but many cases may remain undiagnosed. The inability to diagnose may be due to a unique combination of symptoms or an overlap of conditions, or to the symptoms being atypical of a known disorder, or to the disorder being extremely rare. It is possible that a person with a particular syndrome might not display every single one of the signs and/or symptoms that compose/define a syndrome.
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Positive and negative Sensory symptoms can also be described as positive symptoms, or as negative symptoms depending on whether the symptom is abnormally present such as tingling or itchiness, or abnormally absent such as loss of smell. The following terms are used for negative symptoms – hypoesthesia is a partial loss of sensitivity to moderate stimuli, such as pressure, touch, warmth, cold. Anesthesia is the complete loss of sensitivity to stronger stimuli, such as pinprick. Hypoalgesia (analgesia) is loss of sensation to painful stimuli. Symptoms are also grouped in to negative and positive for some mental disorders such as schizophrenia. Positive symptoms are those that are present in the disorder and are not normally experienced by most individuals and reflects an excess or distortion of normal functions; examples are hallucinations, delusions, and bizarre behavior. Negative symptoms are functions that are normally found but that are diminished or absent, such as apathy and anhedonia. Dynamic and static Dynamic symptoms are capable of change depending on circumstance, whereas static symptoms are fixed or unchanging regardless of circumstance. For example, the symptoms of exercise intolerance are dynamic as they are brought on by exercise, but alleviate during rest. Fixed muscle weakness is a static symptom as the muscle will be weak regardless of exercise or rest. A majority of patients with metabolic myopathies have dynamic rather than static findings, typically experiencing exercise intolerance, muscle pain, and cramps with exercise rather than fixed weakness. Those with the metabolic myopathy of McArdle's disease (GSD-V) and some individuals with phosphoglucomutase deficiency (CDG1T/GSD-XIV), initially experience exercise intolerance during mild-moderate aerobic exercise, but the symptoms alleviate after 6–10 minutes in what is known as "second wind". Neuropsychiatric Neuropsychiatric symptoms are present in many degenerative disorders including dementia, and Parkinson's disease. Symptoms commonly include apathy, anxiety, and depression. Neurological and psychiatric symptoms are also present in some genetic disorders such as Wilson's disease. Symptoms of executive dysfunction are often found in many disorders including schizophrenia, and ADHD. Radiologic Radiologic signs are abnormal medical findings on imaging scanning. These include the Mickey Mouse sign and the Golden S sign. When using imaging to find the cause of a complaint, another unrelated finding may be found known as an incidental finding.
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Cardinal Cardinal signs and symptoms are those that may be diagnostic, and pathognomonic – of a certainty of diagnosis. Inflammation for example has a recognised group of cardinal signs and symptoms, as does exacerbations of chronic bronchitis, and Parkinson's disease. In contrast to a pathognomonic cardinal sign, the absence of a sign or symptom can often rule out a condition. This is known by the Latin term sine qua non. For example, the absence of known genetic mutations specific for a hereditary disease would rule out that disease. Another example is where the vaginal pH is less than 4.5, a diagnosis of bacterial vaginosis would be excluded. Reflexes A reflex is an automatic response in the body to a stimulus. Its absence, reduced (hypoactive), or exaggerated (hyperactive) response can be a sign of damage to the central nervous system or peripheral nervous system. In the patellar reflex (knee-jerk) for example, its reduction or absence is known as Westphal's sign and may indicate damage to lower motor neurons. When the response is exaggerated damage to the upper motor neurons may be indicated. Facies A number of medical conditions are associated with a distinctive facial expression or appearance known as a facies. An example is elfin facies which has facial features like those of the elf, and this may be associated with Williams syndrome, or Donohue syndrome. The most well-known facies is probably the Hippocratic facies that is seen on a person as they near death. Anamnestic signs Anamnestic signs (from anamnēstikós, ἀναμνηστικός, "able to recall to mind") are signs that indicate a past condition, for example paralysis in an arm may indicate a past stroke. Asymptomatic Some diseases including cancers, and infections may be present but show no signs or symptoms and these are known as asymptomatic. A gallstone may be asymptomatic and only discovered as an incidental finding. Easily spreadable viral infections such as COVID-19 may be asymptomatic but may still be transmissible. History
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Symptomatology A symptom (from Greek σύμπτωμα, "accident, misfortune, that which befalls", from συμπίπτω, "I befall", from συν- "together, with" and πίπτω, "I fall") is a departure from normal function or feeling. Symptomatology (also called semiology) is a branch of medicine dealing with the signs and symptoms of a disease. This study also includes the indications of a disease. It was first described as semiotics by Henry Stubbe in 1670 a term now used for the study of sign communication. Prior to the nineteenth century there was little difference in the powers of observation between physician and patient. Most medical practice was conducted as a co-operative interaction between the physician and patient; this was gradually replaced by a "monolithic consensus of opinion imposed from within the community of medical investigators". Whilst each noticed much the same things, the physician had a more informed interpretation of those things: "the physicians knew what the findings meant and the layman did not". Development of medical testing
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A number of advances introduced mostly in the 19th century, allowed for more objective assessment by the physician in search of a diagnosis, and less need of input from the patient. During the 20th century the introduction of a wide range of imaging techniques and other testing methods such as genetic testing, clinical chemistry tests, molecular diagnostics and pathogenomics have made a huge impact on diagnostic capability. In 1761 the percussion technique for diagnosing respiratory conditions was discovered by Leopold Auenbrugger. This method of tapping body cavities to note any abnormal sounds had already been in practice for a long time in cardiology. Percussion of the thorax became more widely known after 1808 with the translation of Auenbrugger's work from Latin into French by Jean-Nicolas Corvisart. In 1819 the introduction of the stethoscope by René Laennec began to replace the centuries-old technique of immediate auscultation – listening to the heart by placing the ear directly on the chest, with mediate auscultation using the stethoscope to listen to the sounds of the heart and respiratory tract. Laennec's publication was translated into English, 1824, by John Forbes. The 1846 introduction by surgeon John Hutchinson (1811–1861) of the spirometer, an apparatus for assessing the mechanical properties of the lungs via measurements of forced exhalation and forced inhalation. (The recorded lung volumes and air flow rates are used to distinguish between restrictive disease (in which the lung volumes are decreased: e.g., cystic fibrosis) and obstructive diseases (in which the lung volume is normal but the air flow rate is impeded; e.g., emphysema).) The 1851 invention by Hermann von Helmholtz (1821–1894) of the ophthalmoscope, which allowed physicians to examine the inside of the human eye. The () immediate widespread clinical use of Sir Thomas Clifford Allbutt's (1836–1925) six-inch (rather than twelve-inch) pocket clinical thermometer, which he had devised in 1867. The 1882 introduction of bacterial cultures by Robert Koch, initially for tuberculosis, being the first laboratory test to confirm bacterial infections. The 1895 clinical use of X-rays which began almost immediately after they had been discovered that year by Wilhelm Conrad Röntgen (1845–1923).
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The 1896 introduction of the sphygmomanometer, designed by Scipione Riva-Rocci (1863–1937), to measure blood pressure.
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Diagnosis The recognition of signs, and noting of symptoms may lead to a diagnosis. Otherwise a physical examination may be carried out, and a medical history taken. Further diagnostic medical tests such as blood tests, scans, and biopsies, may be needed. An X-ray for example would soon be diagnostic of a suspected bone fracture. A noted significance detected during an examination or from a medical test may be known as a medical finding. Examples Ascites (build-up of fluid in the abdomen) Nail clubbing (deformed nails) Cough Death rattle (last moments of life) Hemoptysis (blood-stained sputum) Jaundice Organomegaly an enlarged organ such as the liver (hepatomegaly) Palmar erythema (reddening of hands) Hypersalivation excessive (saliva) Unintentional weight loss
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The sun bear (Helarctos malayanus) is a bear species in the family Ursidae found in the tropical forests of Southeast Asia. It is the only species in the genus Helarctos and the smallest bear species, standing nearly at the shoulder and weighing . It is stockily built, with large paws, strongly curved claws, small, rounded ears and a short snout. The fur is generally short and jet black, but can vary from grey to red. The sun bear gets its name from its characteristic orange to cream-coloured chest patch. Its unique morphology—inward-turned front feet, flattened chest, powerful forelimbs with large claws—suggests adaptations for climbing. The most arboreal (tree-living) of all bears, the sun bear is an excellent climber and sunbathes or sleeps in trees above the ground. It is mainly active during the day, though nocturnality might be more common in areas frequented by humans. Sun bears tend to remain solitary, but sometimes occur in twos (such as a mother and her cub). They do not seem to hibernate, possibly because food resources are available the whole year throughout the range. Being omnivores, sun bears' diet includes ants, bees, beetles, honey, termites, and plant material such as seeds and several kinds of fruits; vertebrates such as birds and deer are also eaten occasionally. They breed throughout the year; individuals become sexually mature at two to four years of age. Litters comprise one or two cubs that remain with their mother for around three years. The range of the sun bear is bounded by northeastern India to the north then south to southeast through Bangladesh, Cambodia, Myanmar, Laos, Thailand, and Vietnam in mainland Asia to Brunei, Indonesia, and Malaysia to the south. These bears are threatened by heavy deforestation and illegal hunting for food and the wildlife trade; they are also harmed in conflicts with humans when they enter farmlands, plantations, and orchards. The global population is estimated to have declined by 35% since the 1990s. The IUCN has listed this species as vulnerable. Etymology The generic name Helarctos comes from two Greek words: (, related to the sun) and (, bear). Another name is honey bear, in Malay and Indonesian, in reference to its habit of feeding on honey from honeycombs.
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Taxonomy and phylogeny The scientific name Ursus malayanus was proposed by Stamford Raffles in 1821 when he first described a sun bear from Sumatra. In 1825, Thomas Horsfield placed the species in a genus of its own, Helarctos, when describing a sun bear from Borneo. Subspecies and distribution H. annamiticus, described by Pierre Marie Heude in 1901 from Annam, is not considered a distinct species, but is subordinated as a junior synonym to H. m. malayanus. In 1906, Richard Lydekker proposed another subspecies by the name H. m. wardii for a sun bear skull, noting its similarities to a skull from Tibet with a thicker coat, but the Tibetan specimen was later found to be an Asian black bear (Ursus thibetanus). Genetic differences between the two subspecies are obscure and some authorities consider the species monotypic. Phylogeny The phylogenetic relationships among ursid species have remained ambiguous over the years. Noting the production of fertile hybrids between sun bears and sloth bears (Melursus ursinus), it was proposed that Helarctos be treated as a synonym of Melursus. However, studies differed on whether the two species were closely related. The brown bear/polar bear genetic lineage was estimated to have genetically diverged from the two black bears/sun bear lineage around (mya); the sun bear appears to have diverged from the two black bears between 6.26 and 5.09 mya. and 5.89–3.51 mya. Nuclear gene sequencing of bear species revealed that the sloth bear and the sun bear were the first Ursinae bears that radiated and are not included in the monophyletic Ursus group; moreover, all relationships between the bears were well resolved. Characteristics
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The sun bear is named so for its characteristic orange- to cream-coloured, crescent-like chest patch. It is the smallest of all bear species. It is stockily built, with large paws, strongly curved claws, small rounded ears and a short snout. The head-and-body length is between , and the shoulder height is nearly . Adults weigh . The snout is grey, silver, or orange. The fur is generally jet black, but can vary from grey to red. The hair is silky and fine, and is the shortest of all bear species, suiting its hot tropical habitat. The characteristic chest patch, typically U-shaped, but sometimes circular or spotlike, varies from orange or ochre-yellow to buff or cream, or even white. Some individuals may even lack the patch. Sun bears can expose the patch while standing on their hind feet as a threat display against enemies. Infants are greyish black with a pale brown or white snout and the chest patch is dirty white; the coat of older juveniles may be dark brown. The underfur is particularly thick and black in adults, while the guard hairs are lighter. Two whorls occur on the shoulders, from whence the hair radiates in all directions. A crest is seen on the sides of the neck and a whorl occurs in the centre of the breast patch. The edges of the paws are tan or brown, and the soles are fur-less, which possibly is an adaptation for climbing trees. The claws are sickle-shaped; the front claws are long and heavy. The tail is long. The sympatric Asian black bear has cream-coloured chest markings of a similar shape as those of sun bears and different claw markings.
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During feeding, the sun bear can extend its exceptionally long tongue to extract insects and honey. The teeth are very large, especially the canines, and the bite force quotient is high relative to its body size for reasons not well understood; a possible explanation could be its frequent opening of tropical hardwood trees with its powerful jaws and claws in pursuit of insects, larvae, or honey. The bite force is high for its size: a 50 Kg sun bear bites with a maximum force of 1907.3–2020.6 Newtons on the rear molar. The head is large, broad and heavy in proportion to the body, but the ears are proportionately smaller; the palate is wide in proportion to the skull. The overall unique morphology of this bear, such as its inward-turned front feet, flattened chest, and powerful fore limbs with large claws, indicates adaptations for extensive climbing. Ecology and behaviour Sun bears lead the most arboreal (tree-living) lifestyle among all bears. They are mainly active during the day, although nocturnality might be more common in areas frequented by humans. The sun bear is an excellent climber; it sunbathes or sleeps in trees above the ground. Bedding sites consist mainly of fallen hollow logs, but they also rest in standing trees with cavities, in cavities underneath fallen logs or tree roots, and in tree branches high above the ground. It is also an efficient swimmer. Sun bears are noted for their intelligence; a captive bear observed sugar being stored in a cupboard locked by a key, and later used its claw to open the lock. A study published in 2019 described skillful mimicry of facial expressions by sun bears, with precision comparable to that seen in gorillas and humans.
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Sun bears are shy and reclusive animals, and usually do not attack humans unless provoked to do so, or if they are injured or with their cubs; their timid nature led these bears to be often tamed and kept as pets in the past. Other sources, though, state that sun bears are known as very fierce animals when surprised in the forest. They are typically solitary but are sometimes seen in pairs (such as mothers and cubs). Sun bears stand on their hind feet for a broader view of their surroundings or smell far-off objects; they try to intimidate their enemies by displaying their chest patch if threatened. Vocalisations include grunts and snuffles while foraging for insects, and roars similar to those of a male orangutan during the breeding season; less commonly, they may give out short barks (like a rhinoceros) when they are surprised. Sun bears do not seem to hibernate, possibly because food resources are available the whole year throughout the range. They occupy home ranges of varying sizes in different areas, ranging from in Borneo and peninsular Malaysia; and in Ulu Segama Forest Reserve in Sabah. Tigers are their major predators; dholes and leopards have also been recorded preying on sun bears, but cases are relatively few. In one incident, a tiger-sun bear interaction resulted in a prolonged altercation and in the death of both animals. In another incident, a wild female sun bear was swallowed by a large reticulated python in East Kalimantan. Diet
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Sun bears are omnivores and feed on a broad variety of items, such as ants, bees, beetles, honey, termites, and plant material such as seeds and several kinds of fruits. Vertebrates such as birds, deer, eggs, and reptiles may be eaten occasionally. They forage mostly at night. Sun bears tear open hollow trees with their long, sharp claws and teeth in search of wild bees and honey. They also break termite mounds and quickly lick and suck the contents, holding pieces of the broken mound with their front paws. They consume figs in large amounts and eat them whole. In a study in the forests of Kalimantan, the fruits of Moraceae, Burseraceae, and Myrtaceae species made up more than 50% of the fruit diet; in times of fruit scarcity, sun bears switched to a more insectivorous diet. A study in Central Borneo revealed that sun bears play an important role in the seed dispersal of Canarium pilosum (a tree in the family Burseraceae). Sun bears eat the centre of coconut palms, and crush oil-rich seeds such as acorns. Oil palms are nutritious but not enough for subsistence. Reproduction Sun bears are polyoestrous; births occur throughout the year. Oestrus lasts five to seven days. Sun bears become sexually mature at two to four years of age. Reported lengths for pregnancies vary from 95 to 240 days; pregnancy tends to be longer in zoos in temperate climate possibly due to delay in implantation or fertilisation. Births occur inside hollow tree cavities. A litter typically comprises one or two cubs weighing around each. Cubs are born deaf with eyes closed. The eyes open at nearly 25 days, but they remain blind till 50 days after birth; the sense of hearing improves over the first 50 days. Cubs younger than two months are dependent on external stimulation for defecation. Cubs are kept on buttress roots at the base of trees until they learn how to walk and climb properly. Mothers protect their cubs aggressively. Offspring remain with their mother for nearly the first three years of their lives. Lifespan in captivity is generally over 20 years; one individual has lived for 34 years. Distribution and habitat
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The sun bear is native to the tropical forests of Southeast Asia; its range is bound by northeastern India to the north and extends south to Bangladesh, Myanmar, Thailand, Cambodia, Laos, and Vietnam to Brunei, Indonesia, and Malaysia to the south. Its presence in China was confirmed in 2017 when it was sighted in Yingjiang County of Yunnan Province. It is extinct in Singapore. These bears dwell primarily in two main types of forests throughout their range - deciduous and seasonally evergreen forests to the north of the Isthmus of Kra, and nonseasonal evergreen forests in Indonesia and Malaysia. They are typically found at low altitudes, such as below in western Thailand and peninsular Malaysia, but this varies widely throughout the range; in India, larger numbers have been recorded at elevations up to than in low-lying areas, probably due to habitat loss at ground level. They occur in montane areas in northeast India, but may not extend farther north into the unfavourable and colder Himalayan region; their distribution might be restricted to the northwest due to competition with sloth bears. The sun bear is sympatric with the Asian black bear throughout the remaining areas in the mainland range featuring a mix of seasonal forest types, with monthly rainfall below for a long spell of 3–7 months. In mountainous areas, Asian black bears are more common than sun bears, probably due to scarcity of invertebrates on which to feed. The major habitats in southern Thailand and peninsular Malaysia are moist evergreen forests, with more or less unvarying climate and heavy rainfall throughout the year, and low-lying or montane dipterocarp forests. Mangroves may be inhabited, but usually only when they are close to preferred habitat types. The sun bear tends to avoid heavily logged forests and areas close to human settlement. However, sun bears have been seen in farmlands, plantations and orchards, where they may be considered vermin. A survey in Lower Kinabatangan Segama Wetlands showed that sun bears were feared but were not common in oil palm plantations; Bornean bearded pigs, elephants and macaques were far more damaging to crops. Sun bears have been reported preying on poultry and livestock.
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Fossil remains suggest its occurrence farther north during the Pleistocene; it may have occurred as far south as Java in the middle to Late Pleistocene. Fossils also known from the Middle Pleistocene of Thailand along with Stegodon, gaur, wild water buffalo, and other living and extinct mammals. Today, it has been eliminated from the majority of its erstwhile range, especially in Thailand; populations are declining in most of the range countries. It disappeared from Singapore during the 1800s and 1900s, possibly due to extensive deforestation. Sun bear populations appear to decrease in size northward from Sundaland, and numbers are especially low in the northern and western extremes of the range. This has possibly been the case since prehistoric times, and is not a result of human interference. The population density varies from in Khao Yai National Park to in the Harapan Rainforest in southern Sumatra. Threats
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According to the IUCN Bear Specialist Group, sun bear populations have fallen by an estimated 35% since the 1990s. Numbers are especially low in Bangladesh and China, and populations in Vietnam are feared to decline severely by 50–80% in the next 30 years. Habitat fragmentation is on the rise, particularly in Borneo, Sumatra, and some areas of the mainland range. Heavy deforestation (due to agriculture, logging, and forest fires) and hunting for wildlife trade are severe threats throughout the range; human-bear conflicts are a relatively minor threat. Compared to other continents, Southeast Asia has undergone severe depletion in forest cover over the past few decades (by almost 12% between 1990 and 2010); this has resulted in substantial habitat loss for forest-dependent species such as sun bears. A 2007 study in East Borneo recorded severe loss of habitat and food resources due to droughts and forest fires brought about by the El Niño. With lack of research in predation, sources have documented very few predation events. In the island of Borneo sun bears were found to be hunted by python in their most vulnerable state. Pythons are successfully able to attack by taking advantage of the nighttime when the sun bears are sleep or nursing their cub. In Southeast Asia, a male leopard (Panthera pardus) has been photographed with a sun bear cub being held by the throat. This reported case has been reported to be the second confirmed predator as of 2019.
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During surveys in Kalimantan between 1994 and 1997, interviewees admitted to hunting sun bears and indicated that sun bear meat is eaten by indigenous people in several areas there. Studies have found evidence of pet trade and sale of sun bear parts such as gall bladders in traditional Chinese medicine (TCM) shops in Sabah and Sarawak. In 2018 and 2019, 128 TCM outlets in 24 locations across Sabah and Sarawak were surveyed and bear parts and derivatives were recorded for sale in 25% of the outlets surveyed, many of which would have been derived from locally sourced sun bears. Sun bears were killed by shooting or administering poison to protect coconut and snakefruit plantations in east Kalimantan. A report published by TRAFFIC in 2011 showed that sun bears, along with Asian black bears and brown bears, are specifically targeted for the bear bile trade in Southeast Asia, and are kept in bear farms in Laos, Vietnam, and Myanmar. Poaching is common in several countries in the region. Hunting pressure is rising even in some protected areas; in the Nam Ha National Protected Area in Laos, hunter snares have been found that specifically target bears. A study in Nagaland (northeastern India) recorded a sparse distribution of sun bears in the Fakim and Ntangki National Parks, and reported extensive illegal hunting for food and trade in bear parts. Protective laws have shown little success in controlling these threats, especially due to poor execution and high potential for gains by the trade.
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Conservation measures The sun bear is listed as vulnerable on the IUCN Red List, and is included in CITES Appendix I. With the exception of Sarawak (Malaysia) and Cambodia, the sun bear is legally protected from hunting in its whole range. A 2014 report documented rampant poaching and trade in sun bear parts in Sarawak, more than anywhere else in Malaysia; the researchers recommended stricter legislations in the state to protect local sun bears. The Bornean Sun Bear Conservation Centre, founded by Wong Siew Te in Sabah (Malaysia) in 2008, aims to work for the welfare of sun bears rescued from poor conditions in captivity and spread awareness about their conservation. The Malayan sun bears are part of an international captive-breeding program and a species survival plan under the Association of Zoos and Aquariums since late 1994. Since that same year, the European breed registry for sun bears is kept in the Cologne Zoological Garden, Germany.
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Membrane potential (also transmembrane potential or membrane voltage) is the difference in electric potential between the interior and the exterior of a biological cell. It equals the interior potential minus the exterior potential. This is the energy (i.e. work) per charge which is required to move a (very small) positive charge at constant velocity across the cell membrane from the exterior to the interior. (If the charge is allowed to change velocity, the change of kinetic energy and production of radiation must be taken into account.) Typical values of membrane potential, normally given in units of milli volts and denoted as mV, range from –80 mV to –40 mV. For such typical negative membrane potentials, positive work is required to move a positive charge from the interior to the exterior. However, thermal kinetic energy allows ions to overcome the potential difference. For a selectively permeable membrane, this permits a net flow against the gradient. This is a kind of osmosis. Description All animal cells are surrounded by a membrane composed of a lipid bilayer with proteins embedded in it. The membrane serves as both an insulator and a diffusion barrier to the movement of ions. Transmembrane proteins, also known as ion transporter or ion pump proteins, actively push ions across the membrane and establish concentration gradients across the membrane, and ion channels allow ions to move across the membrane down those concentration gradients. Ion pumps and ion channels are electrically equivalent to a set of batteries and resistors inserted in the membrane, and therefore create a voltage between the two sides of the membrane. All plasma membranes have an electrical potential across them, with the inside usually negative with respect to the outside. The membrane potential has two basic functions. First, it allows a cell to function as a battery, providing power to operate a variety of "molecular devices" embedded in the membrane. Second, in electrically excitable cells such as neurons and muscle cells, it is used for transmitting signals between different parts of a cell. Signals in neurons and muscle cells Signals are generated in excitable cells by opening or closing of ion channels at one point in the membrane, producing a local change in the membrane potential. This change in the electric field can be quickly sensed by either adjacent or more distant ion channels in the membrane. Those ion channels can then open or close as a result of the potential change, reproducing the signal.
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In non-excitable cells, and in excitable cells in their baseline states, the membrane potential is held at a relatively stable value, called the resting potential. For neurons, resting potential is defined as ranging from –80 to –70 millivolts; that is, the interior of a cell has a negative baseline voltage of a bit less than one-tenth of a volt. The opening and closing of ion channels can induce a departure from the resting potential. This is called a depolarization if the interior voltage becomes less negative (say from –70 mV to –60 mV), or a hyperpolarization if the interior voltage becomes more negative (say from –70 mV to –80 mV). In excitable cells, a sufficiently large depolarization can evoke an action potential, in which the membrane potential changes rapidly and significantly for a short time (on the order of 1 to 100 milliseconds), often reversing its polarity. Action potentials are generated by the activation of certain voltage-gated ion channels. In neurons, the factors that influence the membrane potential are diverse. They include numerous types of ion channels, some of which are chemically gated and some of which are voltage-gated. Because voltage-gated ion channels are controlled by the membrane potential, while the membrane potential itself is influenced by these same ion channels, feedback loops that allow for complex temporal dynamics arise, including oscillations and regenerative events such as action potentials. Ion concentration gradients Differences in the concentrations of ions on opposite sides of a cellular membrane lead to a voltage called the membrane potential. Many ions have a concentration gradient across the membrane, including potassium (K+), which is at a high concentration inside and a low concentration outside the membrane. Sodium (Na+) and chloride (Cl−) ions are at high concentrations in the extracellular region, and low concentrations in the intracellular regions. These concentration gradients provide the potential energy to drive the formation of the membrane potential. This voltage is established when the membrane has permeability to one or more ions. In the simplest case, illustrated in the top diagram ("Ion concentration gradients"), if the membrane is selectively permeable to potassium, these positively charged ions can diffuse down the concentration gradient to the outside of the cell, leaving behind uncompensated negative charges. This separation of charges is what causes the membrane potential.
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The system as a whole is electro-neutral. The uncompensated positive charges outside the cell, and the uncompensated negative charges inside the cell, physically line up on the membrane surface and attract each other across the lipid bilayer. Thus, the membrane potential is physically located only in the immediate vicinity of the membrane. It is the separation of these charges across the membrane that is the basis of the membrane voltage. The top diagram is only an approximation of the ionic contributions to the membrane potential. Other ions including sodium, chloride, calcium, and others play a more minor role, even though they have strong concentration gradients, because they have more limited permeability than potassium. Physical basis The membrane potential in a cell derives ultimately from two factors: electrical force and diffusion. Electrical force arises from the mutual attraction between particles with opposite electrical charges (positive and negative) and the mutual repulsion between particles with the same type of charge (both positive or both negative). Diffusion arises from the statistical tendency of particles to redistribute from regions where they are highly concentrated to regions where the concentration is low. Voltage Voltage, which is synonymous with difference in electrical potential, is the ability to drive an electric current across a resistance. Indeed, the simplest definition of a voltage is given by Ohm's law: V=IR, where V is voltage, I is current and R is resistance. If a voltage source such as a battery is placed in an electrical circuit, the higher the voltage of the source the greater the amount of current that it will drive across the available resistance. The functional significance of voltage lies only in potential differences between two points in a circuit. The idea of a voltage at a single point is meaningless. It is conventional in electronics to assign a voltage of zero to some arbitrarily chosen element of the circuit, and then assign voltages for other elements measured relative to that zero point. There is no significance in which element is chosen as the zero point—the function of a circuit depends only on the differences not on voltages per se. However, in most cases and by convention, the zero level is most often assigned to the portion of a circuit that is in contact with ground.
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The same principle applies to voltage in cell biology. In electrically active tissue, the potential difference between any two points can be measured by inserting an electrode at each point, for example one inside and one outside the cell, and connecting both electrodes to the leads of what is in essence a specialized voltmeter. By convention, the zero potential value is assigned to the outside of the cell and the sign of the potential difference between the outside and the inside is determined by the potential of the inside relative to the outside zero. In mathematical terms, the definition of voltage begins with the concept of an electric field , a vector field assigning a magnitude and direction to each point in space. In many situations, the electric field is a conservative field, which means that it can be expressed as the gradient of a scalar function , that is, . This scalar field is referred to as the voltage distribution. The definition allows for an arbitrary constant of integration—this is why absolute values of voltage are not meaningful. In general, electric fields can be treated as conservative only if magnetic fields do not significantly influence them, but this condition usually applies well to biological tissue. Because the electric field is the gradient of the voltage distribution, rapid changes in voltage within a small region imply a strong electric field; on the converse, if the voltage remains approximately the same over a large region, the electric fields in that region must be weak. A strong electric field, equivalent to a strong voltage gradient, implies that a strong force is exerted on any charged particles that lie within the region. Ions and the forces driving their motion Electrical signals within biological organisms are, in general, driven by ions. The most important cations for the action potential are sodium (Na+) and potassium (K+). Both of these are monovalent cations that carry a single positive charge. Action potentials can also involve calcium (Ca2+), which is a divalent cation that carries a double positive charge. The chloride anion (Cl−) plays a major role in the action potentials of some algae, but plays a negligible role in the action potentials of most animals.
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Ions cross the cell membrane under two influences: diffusion and electric fields. A simple example wherein two solutions—A and B—are separated by a porous barrier illustrates that diffusion will ensure that they will eventually mix into equal solutions. This mixing occurs because of the difference in their concentrations. The region with high concentration will diffuse out toward the region with low concentration. To extend the example, let solution A have 30 sodium ions and 30 chloride ions. Also, let solution B have only 20 sodium ions and 20 chloride ions. Assuming the barrier allows both types of ions to travel through it, then a steady state will be reached whereby both solutions have 25 sodium ions and 25 chloride ions. If, however, the porous barrier is selective to which ions are let through, then diffusion alone will not determine the resulting solution. Returning to the previous example, let's now construct a barrier that is permeable only to sodium ions. Now, only sodium is allowed to diffuse cross the barrier from its higher concentration in solution A to the lower concentration in solution B. This will result in a greater accumulation of sodium ions than chloride ions in solution B and a lesser number of sodium ions than chloride ions in solution A. This means that there is a net positive charge in solution B from the higher concentration of positively charged sodium ions than negatively charged chloride ions. Likewise, there is a net negative charge in solution A from the greater concentration of negative chloride ions than positive sodium ions. Since opposite charges attract and like charges repel, the ions are now also influenced by electrical fields as well as forces of diffusion. Therefore, positive sodium ions will be less likely to travel to the now-more-positive B solution and remain in the now-more-negative A solution. The point at which the forces of the electric fields completely counteract the force due to diffusion is called the equilibrium potential. At this point, the net flow of the specific ion (in this case sodium) is zero. Plasma membranes Every cell is enclosed in a plasma membrane, which has the structure of a lipid bilayer with many types of large molecules embedded in it. Because it is made of lipid molecules, the plasma membrane intrinsically has a high electrical resistivity, in other words a low intrinsic permeability to ions. However, some of the molecules embedded in the membrane are capable either of actively transporting ions from one side of the membrane to the other or of providing channels through which they can move.
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In electrical terminology, the plasma membrane functions as a combined resistor and capacitor. Resistance arises from the fact that the membrane impedes the movement of charges across it. Capacitance arises from the fact that the lipid bilayer is so thin that an accumulation of charged particles on one side gives rise to an electrical force that pulls oppositely charged particles toward the other side. The capacitance of the membrane is relatively unaffected by the molecules that are embedded in it, so it has a more or less invariant value estimated at 2 μF/cm2 (the total capacitance of a patch of membrane is proportional to its area). The conductance of a pure lipid bilayer is so low, on the other hand, that in biological situations it is always dominated by the conductance of alternative pathways provided by embedded molecules. Thus, the capacitance of the membrane is more or less fixed, but the resistance is highly variable. The thickness of a plasma membrane is estimated to be about 7-8 nanometers. Because the membrane is so thin, it does not take a very large transmembrane voltage to create a strong electric field within it. Typical membrane potentials in animal cells are on the order of 100 millivolts (that is, one tenth of a volt), but calculations show that this generates an electric field close to the maximum that the membrane can sustain—it has been calculated that a voltage difference much larger than 200 millivolts could cause dielectric breakdown, that is, arcing across the membrane. Facilitated diffusion and transport The resistance of a pure lipid bilayer to the passage of ions across it is very high, but structures embedded in the membrane can greatly enhance ion movement, either actively or passively, via mechanisms called facilitated transport and facilitated diffusion. The two types of structure that play the largest roles are ion channels and ion pumps, both usually formed from assemblages of protein molecules. Ion channels provide passageways through which ions can move. In most cases, an ion channel is permeable only to specific types of ions (for example, sodium and potassium but not chloride or calcium), and sometimes the permeability varies depending on the direction of ion movement. Ion pumps, also known as ion transporters or carrier proteins, actively transport specific types of ions from one side of the membrane to the other, sometimes using energy derived from metabolic processes to do so. Ion pumps
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Ion pumps are integral membrane proteins that carry out active transport, i.e., use cellular energy (ATP) to "pump" the ions against their concentration gradient. Such ion pumps take in ions from one side of the membrane (decreasing its concentration there) and release them on the other side (increasing its concentration there). The ion pump most relevant to the action potential is the sodium–potassium pump, which transports three sodium ions out of the cell and two potassium ions in. As a consequence, the concentration of potassium ions K+ inside the neuron is roughly 30-fold larger than the outside concentration, whereas the sodium concentration outside is roughly five-fold larger than inside. In a similar manner, other ions have different concentrations inside and outside the neuron, such as calcium, chloride and magnesium. If the numbers of each type of ion were equal, the sodium–potassium pump would be electrically neutral, but, because of the three-for-two exchange, it gives a net movement of one positive charge from intracellular to extracellular for each cycle, thereby contributing to a positive voltage difference. The pump has three effects: (1) it makes the sodium concentration high in the extracellular space and low in the intracellular space; (2) it makes the potassium concentration high in the intracellular space and low in the extracellular space; (3) it gives the intracellular space a negative voltage with respect to the extracellular space. The sodium-potassium pump is relatively slow in operation. If a cell were initialized with equal concentrations of sodium and potassium everywhere, it would take hours for the pump to establish equilibrium. The pump operates constantly, but becomes progressively less efficient as the concentrations of sodium and potassium available for pumping are reduced. Ion pumps influence the action potential only by establishing the relative ratio of intracellular and extracellular ion concentrations. The action potential involves mainly the opening and closing of ion channels not ion pumps. If the ion pumps are turned off by removing their energy source, or by adding an inhibitor such as ouabain, the axon can still fire hundreds of thousands of action potentials before their amplitudes begin to decay significantly. In particular, ion pumps play no significant role in the repolarization of the membrane after an action potential.
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Another functionally important ion pump is the sodium-calcium exchanger. This pump operates in a conceptually similar way to the sodium-potassium pump, except that in each cycle it exchanges three Na+ from the extracellular space for one Ca++ from the intracellular space. Because the net flow of charge is inward, this pump runs "downhill", in effect, and therefore does not require any energy source except the membrane voltage. Its most important effect is to pump calcium outward—it also allows an inward flow of sodium, thereby counteracting the sodium-potassium pump, but, because overall sodium and potassium concentrations are much higher than calcium concentrations, this effect is relatively unimportant. The net result of the sodium-calcium exchanger is that in the resting state, intracellular calcium concentrations become very low. Ion channels Ion channels are integral membrane proteins with a pore through which ions can travel between extracellular space and cell interior. Most channels are specific (selective) for one ion; for example, most potassium channels are characterized by 1000:1 selectivity ratio for potassium over sodium, though potassium and sodium ions have the same charge and differ only slightly in their radius. The channel pore is typically so small that ions must pass through it in single-file order. Channel pores can be either open or closed for ion passage, although a number of channels demonstrate various sub-conductance levels. When a channel is open, ions permeate through the channel pore down the transmembrane concentration gradient for that particular ion. Rate of ionic flow through the channel, i.e. single-channel current amplitude, is determined by the maximum channel conductance and electrochemical driving force for that ion, which is the difference between the instantaneous value of the membrane potential and the value of the reversal potential. A channel may have several different states (corresponding to different conformations of the protein), but each such state is either open or closed. In general, closed states correspond either to a contraction of the pore—making it impassable to the ion—or to a separate part of the protein, stoppering the pore. For example, the voltage-dependent sodium channel undergoes inactivation, in which a portion of the protein swings into the pore, sealing it. This inactivation shuts off the sodium current and plays a critical role in the action potential.
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Ion channels can be classified by how they respond to their environment. For example, the ion channels involved in the action potential are voltage-sensitive channels; they open and close in response to the voltage across the membrane. Ligand-gated channels form another important class; these ion channels open and close in response to the binding of a ligand molecule, such as a neurotransmitter. Other ion channels open and close with mechanical forces. Still other ion channels—such as those of sensory neurons—open and close in response to other stimuli, such as light, temperature or pressure. Leakage channels Leakage channels are the simplest type of ion channel, in that their permeability is more or less constant. The types of leakage channels that have the greatest significance in neurons are potassium and chloride channels. Even these are not perfectly constant in their properties: First, most of them are voltage-dependent in the sense that they conduct better in one direction than the other (in other words, they are rectifiers); second, some of them are capable of being shut off by chemical ligands even though they do not require ligands in order to operate. Ligand-gated channels Ligand-gated ion channels are channels whose permeability is greatly increased when some type of chemical ligand binds to the protein structure. Animal cells contain hundreds, if not thousands, of types of these. A large subset function as neurotransmitter receptors—they occur at postsynaptic sites, and the chemical ligand that gates them is released by the presynaptic axon terminal. One example of this type is the AMPA receptor, a receptor for the neurotransmitter glutamate that when activated allows passage of sodium and potassium ions. Another example is the GABAA receptor, a receptor for the neurotransmitter GABA that when activated allows passage of chloride ions. Neurotransmitter receptors are activated by ligands that appear in the extracellular area, but there are other types of ligand-gated channels that are controlled by interactions on the intracellular side. Voltage-dependent channels Voltage-gated ion channels, also known as voltage dependent ion channels, are channels whose permeability is influenced by the membrane potential. They form another very large group, with each member having a particular ion selectivity and a particular voltage dependence. Many are also time-dependent—in other words, they do not respond immediately to a voltage change but only after a delay.
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One of the most important members of this group is a type of voltage-gated sodium channel that underlies action potentials—these are sometimes called Hodgkin-Huxley sodium channels because they were initially characterized by Alan Lloyd Hodgkin and Andrew Huxley in their Nobel Prize-winning studies of the physiology of the action potential. The channel is closed at the resting voltage level, but opens abruptly when the voltage exceeds a certain threshold, allowing a large influx of sodium ions that produces a very rapid change in the membrane potential. Recovery from an action potential is partly dependent on a type of voltage-gated potassium channel that is closed at the resting voltage level but opens as a consequence of the large voltage change produced during the action potential. Reversal potential The reversal potential (or equilibrium potential) of an ion is the value of transmembrane voltage at which diffusive and electrical forces counterbalance, so that there is no net ion flow across the membrane. This means that the transmembrane voltage exactly opposes the force of diffusion of the ion, such that the net current of the ion across the membrane is zero and unchanging. The reversal potential is important because it gives the voltage that acts on channels permeable to that ion—in other words, it gives the voltage that the ion concentration gradient generates when it acts as a battery. The equilibrium potential of a particular ion is usually designated by the notation Eion.The equilibrium potential for any ion can be calculated using the Nernst equation. For example, reversal potential for potassium ions will be as follows: where Eeq,K+= equilibrium potential for potassium, measured in volts R = universal gas constant, equal to 8.314 joules·K−1·mol−1 T = absolute temperature, measured in kelvins (= K = degrees Celsius + 273.15) z = number of elementary charges of the ion in question involved in the reaction F = Faraday constant, equal to 96,485 coulombs·mol−1 or J·V−1·mol−1 [K+]o= extracellular concentration of potassium, measured in mol·m−3 or mmol·l−1 [K+]i= intracellular concentration of potassium
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Even if two different ions have the same charge (i.e., K+ and Na+), they can still have very different equilibrium potentials, provided their outside and/or inside concentrations differ. Take, for example, the equilibrium potentials of potassium and sodium in neurons. The potassium equilibrium potential EK is −84 mV with 5 mM potassium outside and 140 mM inside. On the other hand, the sodium equilibrium potential, ENa, is approximately +66 mV with approximately 12 mM sodium inside and 140 mM outside. Changes to membrane potential during development A neuron's resting membrane potential actually changes during the development of an organism. In order for a neuron to eventually adopt its full adult function, its potential must be tightly regulated during development. As an organism progresses through development the resting membrane potential becomes more negative. Glial cells are also differentiating and proliferating as development progresses in the brain. The addition of these glial cells increases the organism's ability to regulate extracellular potassium. The drop in extracellular potassium can lead to a decrease in membrane potential of 35 mV. Cell excitability Cell excitability is the change in membrane potential that is necessary for cellular responses in various tissues. Cell excitability is a property that is induced during early embriogenesis. Excitability of a cell has also been defined as the ease with which a response may be triggered. The resting and threshold potentials forms the basis of cell excitability and these processes are fundamental for the generation of graded and action potentials.
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The most important regulators of cell excitability are the extracellular electrolyte concentrations (i.e. Na+, K+, Ca2+, Cl−, Mg2+) and associated proteins. Important proteins that regulate cell excitability are voltage-gated ion channels, ion transporters (e.g. Na+/K+-ATPase, magnesium transporters, acid–base transporters), membrane receptors and hyperpolarization-activated cyclic-nucleotide-gated channels. For example, potassium channels and calcium-sensing receptors are important regulators of excitability in neurons, cardiac myocytes and many other excitable cells like astrocytes. Calcium ion is also the most important second messenger in excitable cell signaling. Activation of synaptic receptors initiates long-lasting changes in neuronal excitability. Thyroid, adrenal and other hormones also regulate cell excitability, for example, progesterone and estrogen modulate myometrial smooth muscle cell excitability. Many cell types are considered to have an excitable membrane. Excitable cells are neurons, muscle (cardiac, skeletal, smooth), vascular endothelial cells, pericytes, juxtaglomerular cells, interstitial cells of Cajal, many types of epithelial cells (e.g. beta cells, alpha cells, delta cells, enteroendocrine cells, pulmonary neuroendocrine cells, pinealocytes), glial cells (e.g. astrocytes), mechanoreceptor cells (e.g. hair cells and Merkel cells), chemoreceptor cells (e.g. glomus cells, taste receptors), some plant cells and possibly immune cells. Astrocytes display a form of non-electrical excitability based on intracellular calcium variations related to the expression of several receptors through which they can detect the synaptic signal. In neurons, there are different membrane properties in some portions of the cell, for example, dendritic excitability endows neurons with the capacity for coincidence detection of spatially separated inputs. Equivalent circuit
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Electrophysiologists model the effects of ionic concentration differences, ion channels, and membrane capacitance in terms of an equivalent circuit, which is intended to represent the electrical properties of a small patch of membrane. The equivalent circuit consists of a capacitor in parallel with four pathways each consisting of a battery in series with a variable conductance. The capacitance is determined by the properties of the lipid bilayer, and is taken to be fixed. Each of the four parallel pathways comes from one of the principal ions, sodium, potassium, chloride, and calcium. The voltage of each ionic pathway is determined by the concentrations of the ion on each side of the membrane; see the Reversal potential section above. The conductance of each ionic pathway at any point in time is determined by the states of all the ion channels that are potentially permeable to that ion, including leakage channels, ligand-gated channels, and voltage-gated ion channels. For fixed ion concentrations and fixed values of ion channel conductance, the equivalent circuit can be further reduced, using the Goldman equation as described below, to a circuit containing a capacitance in parallel with a battery and conductance. In electrical terms, this is a type of RC circuit (resistance-capacitance circuit), and its electrical properties are very simple. Starting from any initial state, the current flowing across either the conductance or the capacitance decays with an exponential time course, with a time constant of , where is the capacitance of the membrane patch, and is the net resistance. For realistic situations, the time constant usually lies in the 1—100 millisecond range. In most cases, changes in the conductance of ion channels occur on a faster time scale, so an RC circuit is not a good approximation; however, the differential equation used to model a membrane patch is commonly a modified version of the RC circuit equation.
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Resting potential When the membrane potential of a cell goes for a long period of time without changing significantly, it is referred to as a resting potential or resting voltage. This term is used for the membrane potential of non-excitable cells, but also for the membrane potential of excitable cells in the absence of excitation. In excitable cells, the other possible states are graded membrane potentials (of variable amplitude), and action potentials, which are large, all-or-nothing rises in membrane potential that usually follow a fixed time course. Excitable cells include neurons, muscle cells, and some secretory cells in glands. Even in other types of cells, however, the membrane voltage can undergo changes in response to environmental or intracellular stimuli. For example, depolarization of the plasma membrane appears to be an important step in programmed cell death. The interactions that generate the resting potential are modeled by the Goldman equation. This is similar in form to the Nernst equation shown above, in that it is based on the charges of the ions in question, as well as the difference between their inside and outside concentrations. However, it also takes into consideration the relative permeability of the plasma membrane to each ion in question. The three ions that appear in this equation are potassium (K+), sodium (Na+), and chloride (Cl−). Calcium is omitted, but can be added to deal with situations in which it plays a significant role. Being an anion, the chloride terms are treated differently from the cation terms; the intracellular concentration is in the numerator, and the extracellular concentration in the denominator, which is reversed from the cation terms. Pi stands for the relative permeability of the ion type i.
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In essence, the Goldman formula expresses the membrane potential as a weighted average of the reversal potentials for the individual ion types, weighted by permeability. (Although the membrane potential changes about 100 mV during an action potential, the concentrations of ions inside and outside the cell do not change significantly. They remain close to their respective concentrations when then membrane is at resting potential.) In most animal cells, the permeability to potassium is much higher in the resting state than the permeability to sodium. As a consequence, the resting potential is usually close to the potassium reversal potential. The permeability to chloride can be high enough to be significant, but, unlike the other ions, chloride is not actively pumped, and therefore equilibrates at a reversal potential very close to the resting potential determined by the other ions. Values of resting membrane potential in most animal cells usually vary between the potassium reversal potential (usually around -80 mV) and around -40 mV. The resting potential in excitable cells (capable of producing action potentials) is usually near -60 mV—more depolarized voltages would lead to spontaneous generation of action potentials. Immature or undifferentiated cells show highly variable values of resting voltage, usually significantly more positive than in differentiated cells. In such cells, the resting potential value correlates with the degree of differentiation: undifferentiated cells in some cases may not show any transmembrane voltage difference at all. Maintenance of the resting potential can be metabolically costly for a cell because of its requirement for active pumping of ions to counteract losses due to leakage channels. The cost is highest when the cell function requires an especially depolarized value of membrane voltage. For example, the resting potential in daylight-adapted blowfly (Calliphora vicina) photoreceptors can be as high as -30 mV. This elevated membrane potential allows the cells to respond very rapidly to visual inputs; the cost is that maintenance of the resting potential may consume more than 20% of overall cellular ATP.
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On the other hand, the high resting potential in undifferentiated cells does not necessarily incur a high metabolic cost. This apparent paradox is resolved by examination of the origin of that resting potential. Little-differentiated cells are characterized by extremely high input resistance, which implies that few leakage channels are present at this stage of cell life. As an apparent result, potassium permeability becomes similar to that for sodium ions, which places resting potential in-between the reversal potentials for sodium and potassium as discussed above. The reduced leakage currents also mean there is little need for active pumping in order to compensate, therefore low metabolic cost. Graded potentials As explained above, the potential at any point in a cell's membrane is determined by the ion concentration differences between the intracellular and extracellular areas, and by the permeability of the membrane to each type of ion. The ion concentrations do not normally change very quickly (with the exception of Ca2+, where the baseline intracellular concentration is so low that even a small influx may increase it by orders of magnitude), but the permeabilities of the ions can change in a fraction of a millisecond, as a result of activation of ligand-gated ion channels. The change in membrane potential can be either large or small, depending on how many ion channels are activated and what type they are, and can be either long or short, depending on the lengths of time that the channels remain open. Changes of this type are referred to as graded potentials, in contrast to action potentials, which have a fixed amplitude and time course. As can be derived from the Goldman equation shown above, the effect of increasing the permeability of a membrane to a particular type of ion shifts the membrane potential toward the reversal potential for that ion. Thus, opening Na+ channels shifts the membrane potential toward the Na+ reversal potential, which is usually around +100 mV. Likewise, opening K+ channels shifts the membrane potential toward about –90 mV, and opening Cl− channels shifts it toward about –70 mV (resting potential of most membranes). Thus, Na+ channels shift the membrane potential in a positive direction, K+ channels shift it in a negative direction (except when the membrane is hyperpolarized to a value more negative than the K+ reversal potential), and Cl− channels tend to shift it towards the resting potential.
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Graded membrane potentials are particularly important in neurons, where they are produced by synapses—a temporary change in membrane potential produced by activation of a synapse by a single graded or action potential is called a postsynaptic potential. Neurotransmitters that act to open Na+ channels typically cause the membrane potential to become more positive, while neurotransmitters that activate K+ channels typically cause it to become more negative; those that inhibit these channels tend to have the opposite effect. Whether a postsynaptic potential is considered excitatory or inhibitory depends on the reversal potential for the ions of that current, and the threshold for the cell to fire an action potential (around –50mV). A postsynaptic current with a reversal potential above threshold, such as a typical Na+ current, is considered excitatory. A current with a reversal potential below threshold, such as a typical K+ current, is considered inhibitory. A current with a reversal potential above the resting potential, but below threshold, will not by itself elicit action potentials, but will produce subthreshold membrane potential oscillations. Thus, neurotransmitters that act to open Na+ channels produce excitatory postsynaptic potentials, or EPSPs, whereas neurotransmitters that act to open K+ or Cl− channels typically produce inhibitory postsynaptic potentials, or IPSPs. When multiple types of channels are open within the same time period, their postsynaptic potentials summate (are added together). Other values From the viewpoint of biophysics, the resting membrane potential is merely the membrane potential that results from the membrane permeabilities that predominate when the cell is resting. The above equation of weighted averages always applies, but the following approach may be more easily visualized. At any given moment, there are two factors for an ion that determine how much influence that ion will have over the membrane potential of a cell: That ion's driving force That ion's permeability
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If the driving force is high, then the ion is being "pushed" across the membrane. If the permeability is high, it will be easier for the ion to diffuse across the membrane. Driving force is the net electrical force available to move that ion across the membrane. It is calculated as the difference between the voltage that the ion "wants" to be at (its equilibrium potential) and the actual membrane potential (Em). So, in formal terms, the driving force for an ion = Em - Eion For example, at our earlier calculated resting potential of −73 mV, the driving force on potassium is 7 mV : (−73 mV) − (−80 mV) = 7 mV. The driving force on sodium would be (−73 mV) − (60 mV) = −133 mV. Permeability is a measure of how easily an ion can cross the membrane. It is normally measured as the (electrical) conductance and the unit, siemens, corresponds to 1 C·s−1·V−1, that is one coulomb per second per volt of potential. So, in a resting membrane, while the driving force for potassium is low, its permeability is very high. Sodium has a huge driving force but almost no resting permeability. In this case, potassium carries about 20 times more current than sodium, and thus has 20 times more influence over Em than does sodium. However, consider another case—the peak of the action potential. Here, permeability to Na is high and K permeability is relatively low. Thus, the membrane moves to near ENa and far from EK. The more ions are permeant the more complicated it becomes to predict the membrane potential. However, this can be done using the Goldman-Hodgkin-Katz equation or the weighted means equation. By plugging in the concentration gradients and the permeabilities of the ions at any instant in time, one can determine the membrane potential at that moment. What the GHK equations means is that, at any time, the value of the membrane potential will be a weighted average of the equilibrium potentials of all permeant ions. The "weighting" is the ions relative permeability across the membrane.
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Effects and implications While cells expend energy to transport ions and establish a transmembrane potential, they use this potential in turn to transport other ions and metabolites such as sugar. The transmembrane potential of the mitochondria drives the production of ATP, which is the common currency of biological energy. Cells may draw on the energy they store in the resting potential to drive action potentials or other forms of excitation. These changes in the membrane potential enable communication with other cells (as with action potentials) or initiate changes inside the cell, which happens in an egg when it is fertilized by a sperm. Changes in the dielectric properties of plasma membrane may act as hallmark of underlying conditions such as diabetes and dyslipidemia. In neuronal cells, an action potential begins with a rush of sodium ions into the cell through sodium channels, resulting in depolarization, while recovery involves an outward rush of potassium through potassium channels. Both of these fluxes occur by passive diffusion. A dose of salt may trigger the still-working neurons of a fresh cut of meat into firing, causing muscle spasms.
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Adansonia is a genus made up of eight species of medium-to-large deciduous trees known as baobabs ( or ) or adansonias. They are placed in the family Malvaceae, subfamily Bombacoideae. They are native to Madagascar, mainland Africa, and Australia. The trees have also been introduced to other regions such as Asia. A genomic and ecological analysis has suggested that the genus is Madagascan in origin. The generic name honours Michel Adanson, the French naturalist and explorer who described Adansonia digitata. The baobab is also known as the "upside down tree", a name that originates from its appearance and several myths. They are among the most long-lived of vascular plants and have large flowers that are reproductive for a maximum of 15 hours. The flowers open around dusk, opening so quickly that movement can be detected by the naked eye, and are faded by the next morning. The fruits are large, oval to round and berry-like and hold kidney-shaped seeds in a dry, pulpy matrix. In the early 21st century, baobabs in southern Africa began to die off rapidly from a cause yet to be determined. It is unlikely that disease or pests would be able to kill many trees so rapidly, and some have speculated that the die-off is a result of dehydration. Description Baobabs are long-lived deciduous, small to large trees from tall with broad trunks and compact crowns. Young trees usually have slender, tapering trunks, often with a swollen base. Mature trees have massive trunks that are bottle-shaped or cylindrical and tapered from bottom to top. The trunk is made of fibrous wood arranged in concentric rings, although rings are not always formed annually and so cannot be used to determine the age of individual trees. Tree diameter fluctuates with rainfall so it is thought that water may be stored in the trunk. Baobab trees have two types of shoots—long, green vegetative ones, and stout, woody reproductive ones. Branches can be massive and spread out horizontally from the trunk or are ascending. Adansonia gregorii is generally the smallest of the baobabs, rarely getting to over tall and often with multiple trunks. Both A. rubrostipa and A. madagascariensis are small to large trees, from tall. The other baobabs grow from tall, with diameter trunks. A. digitata, however, often has massive single or multiple trunks of up to diameter.
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