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3596398 | https://en.wikipedia.org/wiki/Polarography | Polarography | Polarography is a type of voltammetry where the working electrode is a dropping mercury electrode (DME) or a static mercury drop electrode (SMDE), which are useful for their wide cathodic ranges and renewable surfaces. It was invented in 1922 by Czechoslovak chemist Jaroslav Heyrovský, for which he won the Nobel prize in 1959. The main advantages of mercury as electrode material are as follows:
1) a large voltage window: ca. from +0.2 V to -1.8 V vs reversible hydrogen electrode (RHE). Hg electrode is particularly well-suited for studying electroreduction reactions.
2) very reproducible electrode surface, since mercury is liquid.
3) very easy cleaning of the electrode surface by making a new drop of mercury from a large Hg pool connected by a glass capillary.
Polarography played a major role as an experimental tool in the advancement of both Analytical Chemistry and Electrochemistry until the 1990s (see figure below), when it was supplanted by other methods that did not require the use of mercury.
Principle of operation
Polarography is an electrochemical voltammetric technique that employs (dropping or static) mercury drop as a working electrode. In its most simple form polarography can be used to determine concentrations of electroactive species in liquids by measuring their mass-transport limiting currents. In such an experiment the potential of the working mercury drop electrode is linearly changed in time, and the electrode current is recorded at a certain time just before the mercury drop dislodges from a glass capillary from where the stream of mercury emerges. A plot of the current vs. potential in a polarography experiment shows the current oscillations corresponding to the drops of Hg falling from the capillary. If the maximum currents of each drop were connected, a sigmoidal shape would result. The limiting current (the plateau on the sigmoid), is called the diffusion-limited current because diffusion is the principal contribution to the flux of the electroactive material at this point of the Hg drop life. More advanced varieties of polarography (see below) produce peaks (which allow for a better resolution of different chemical species) rather than the waves of classical polarography, and improve the detection limits, which in some cases can be as low as 10^-9 M.
Limitations
There are limitations in particular for the classical polarography experiment for quantitative analytical measurements. Because the current is continuously measured during the growth of the Hg drop, there is a substantial contribution from capacitive current. As the Hg flows from the capillary end, there is initially a large increase in the surface area. As a consequence, the initial current is dominated by capacitive effects as charging of the rapidly increasing interface occurs. Toward the end of the drop life, there is little change in the surface area which diminishes the contribution of capacitance changes to the total current. At the same time, any redox process which occurs will result in faradaic current that decays approximately as the square root of time (due to the increasing dimensions of the Nernst diffusion layer). The exponential decay of the capacitive current is much more rapid than the decay of the faradaic current; hence, the faradaic current is proportionally larger at the end of the drop life. Unfortunately, this process is complicated by the continuously changing potential that is applied to the working electrode (the Hg drop) throughout the experiment. Because the potential changes during the drop lifetime (assuming typical experimental parameters of a 2 mV/s scan rate and a 4 s drop time, the potential can change by 8 mV from the beginning to the end of the drop), the charging of the interface (capacitive current) has a continuous contribution to the total current, even at the end of the drop when the surface area is not rapidly changing. As such, the typical signal to noise ratio of a polarographic experiment allows detection limits of only approximately 10−5 or 10−6 M.
Improvements
Dramatically better discrimination against the capacitive current can be obtained using the tast and pulse polarographic techniques. These have been developed with the introduction of analogue and digital electronic potentiostats. The first major improvement was obtained by measuring the current only at the end of each drop lifetime (tast polarography). An even greater enhancement was the introduction of differential pulse polarography. Here, the current is measured before the beginning and before the end of short potential pulses. The latter are superimposed on the linear potential-time-function of the voltammetric scan. Typical amplitudes of these pulses range between 10 and 50 mV, whereas pulse duration is 20 to 50 ms. The difference between both current values is the analytical signal. This technique results in a 100 to 1000-fold improvement of the detection limit, because the capacitive component is effectively subtracted.
Qualitative information
Qualitative information can also be determined from the half-wave potential of the polarogram (the current vs. potential plot in a polarographic experiment). The value of the half-wave potential is related to the standard potential for the redox reaction being studied.
This technique and especially the differential pulse anodic stripping voltammetry (DPASV) method can be used for environmental analysis, and especially for marine study for the characterisation of organic matter and metals interactions.
Quantitative information
The Ilkovic equation is a relation used in polarography relating the diffusion current (Id) and the concentration of the depolarizer (c), which is the substance reduced or oxidized at the dropping mercury electrode. The Ilkovic equation has the form
where:
k is a constant which includes π and the density of mercury, and with the Faraday constant F has been evaluated at 708 for maximal current and 607 for average current
D is the diffusion coefficient of the depolarizer in the medium (cm2/s)
n is the number of electrons exchanged in the electrode reaction, m is the mass flow rate of Hg through the capillary (mg/s)
t is the drop lifetime in seconds,
c is depolarizer concentration in mol/cm3.
The equation is named after the scientist who derived it, the Slovak chemist Dionýz Ilkovič (1907–1980).
| Physical sciences | Electrical methods | Chemistry |
132151 | https://en.wikipedia.org/wiki/Helmet | Helmet | A helmet is a form of protective gear worn to protect the head. More specifically, a helmet complements the skull in protecting the human brain. Ceremonial or symbolic helmets (e.g., a policeman's helmet in the United Kingdom) without protective function are sometimes worn. Soldiers wear combat helmets, often made from Kevlar or other lightweight synthetic fibers.
The word helmet is derived from helm, an Old English word for a protective head covering.
Helmets are used for most sports (e.g., jockeys, American football, ice hockey, cricket, baseball, skiing, hurling and rock climbing); dangerous work activities such as construction, mining, riot police, military aviation, and in transportation (e.g. motorcycle helmets and bicycle helmets). Since the 1990s, most helmets are made from resin or plastic, which may be reinforced with fibers such as aramids.
Designs
Some British gamekeepers during the 18th and 19th centuries wore helmets made of straw bound together with cut bramble. Europeans in the tropics often wore the pith helmet, developed in the mid-19th century and made of pith or cork.
Military applications in the 19th–20th centuries saw a number of leather helmets, particularly among aviators and tank crews in the early 20th century. In the early days of the automobile, some motorists also adopted this style of headgear, and early football helmets were also made of leather. In World War II, American, Soviet, German, Italian and French flight crews wore leather helmets, the German pilots disguising theirs under a beret before disposing of both and switching to cloth caps. The era of the First and Second World Wars also saw a resurgence of metal military helmets, most notably the Brodie helmet and the Stahlhelm.
Modern helmets have a much wider range of applications, including helmets adapted to the specific needs of many athletic pursuits and work environments, and these helmets very often incorporate plastics and other synthetic materials for their light weight and shock absorption capabilities. Some types of synthetic fibers used to make helmets in the 21st century include aramid fibers, such as Kevlar and Twaron. Race car helmets include a head and neck support system that keeps the helmet (and head) attached to the body in severe collisions.
Helmet types
Helmets of many different types have developed over time. Most early helmets had military uses, though some may have had more ceremonial than combat applications.
Two important helmet types to develop in antiquity were the Corinthian helmet and the Roman galea.
During the Middle Ages, many different military helmets and some ceremonial helmets were developed, almost all being metal. Some of the more important medieval developments included the great helm, the bascinet, the frog-mouth helm, and the armet.
The great seal of Owain Glyndŵr (c. 1359 – c. 1415) depicts the prince of Wales & his stallion wearing full armour, they both wear protective headgear with Owain's gold dragon mounted on top. This would have been impractical in battle, so therefore these would have been ceremonial.
In the 19th century, more materials were incorporated, namely leather, felt and pith. The pith helmet and the leather pickelhaube were important 19th century developments. The greatest expansion in the variety of forms and composition of helmets, however, took place in the 20th century, with the development of highly specialized helmets for a multitude of athletic and professional applications, as well as the advent of modern plastics. During World War I, the French army developed the Adrian helmet, the British developed the Brodie helmet, and the Germans produced the Stahlhelm.
The development of hard hats for workplace safety may have been inspired by the helmets of WWI, and they have become a standard type of safety equipment on many construction job sites and industrial locations.
Flight helmets were also developed throughout the 20th century. A multitude of athletic helmets, including football helmets, batting helmets, hockey helmets, cricket helmets, bicycle helmets, ski helmets, motorcycle helmets and racing helmets, were also developed in the 20th century.
Helmets since the mid-20th century have often incorporated lightweight plastics and other synthetic materials, and their use has become highly specialized. Some important recent developments include the French SPECTRA helmet, Spanish MARTE helmet or the American PASGT (commonly called "Kevlar" by U.S. troops) and Advanced Combat Helmet, or ACH.
Heraldry
As the coat of arms was originally designed to distinguish noble combatants on the battlefield or in a tournament, even while covered in armour, it is not surprising that heraldic elements constantly incorporated the shield and the helmet, these often being the most visible parts of a knight's military equipment.
The practice of indicating peerage through the display of barred or grilled helmets first appeared around 1587-1615, and the heraldic convention of displaying helmets of rank in the United Kingdom, which came into vogue around Stuart times, is as follows:
Sovereign: a gold barred-face (tournament) helm placed affronté
Peer's helmet: silver barred-face (tournament) helm placed in profile
Knight's or baronet's helmet: steel helm (earlier jousting helm, later close helm) placed affronté with visor open
Esquire's helmet: steel helm placed in profile with visor closed
Earlier rolls of arms reveal, however, that early heraldic helmets were depicted in a manner faithful to the styles in actual military or tournament use at the time.
Gallery
| Technology | Armour | null |
25299326 | https://en.wikipedia.org/wiki/Chronic%20periodontitis | Chronic periodontitis | Chronic periodontitis is one of the seven categories of periodontitis as defined by the American Academy of Periodontology 1999 classification system. Chronic periodontitis is a common disease of the oral cavity consisting of chronic inflammation of the periodontal tissues that is caused by the accumulation of profuse amounts of dental plaque. Periodontitis initially begins as gingivitis and can progress onto chronic and subsequent aggressive periodontitis according to the 1999 classification.
Diagnosing chronic periodontitis is important in its early stages to prevent severe and irreversible damage to the protective and supportive structures of the tooth. However, due to chronic periodontitis being a painless progressing disease, few patients will seek dental care in the early stages. Mild to moderate chronic periodontitis can be managed by proper mechanical removal of the biofilm and calculus subgingivally. Full and effective oral hygiene and regular 3 monthly periodontal checkups are important for maintaining the stability of the disease.
Chronic periodontitis is prevalent in adults and seniors worldwide. In the US around 35% of adults (30–90 years) are affected. The cumulative effects of alveolar bone loss, attachment loss and pocket formation is more apparent with an increase in age. Age is related to the incidence of periodontal destruction: "...in a well-maintained population who practises oral home care and has regular check-ups, the incidence of incipient periodontal destruction increases with age, the highest rate occurs between 50 and 60 years, and gingival recession is the predominant lesion before 40 years, while periodontal pocketing is the principal mode of destruction between 50 and 60 years of age."
There are a variety of periodontal risk factors which can affect the prevalence, rate, extent and severity of the disease progression. Major risk factors include smoking, lack of oral hygiene with inadequate plaque biofilm control.
There is a slow to moderate rate of disease progression but the patient may have periods of rapid progression ("bursts of destruction"). Chronic periodontitis can be associated with local predisposing factors (e.g. tooth-related or iatrogenic factors). The disease may be modified by and be associated with systemic diseases (e.g. diabetes mellitus, HIV infection) It can also be modified by factors other than systemic disease such as smoking and emotional stress, anxiety and depression. Care should be taken however, when diagnosing a patient who smokes as smoking can alter some of the results of an examination. In smokers, the gingiva are pale and fibrous and tend to bleed less while being probed due to the effect of nicotine on the vasculature by vasoconstricting them. Thus, a lowered response is produced and this explains why incorrect data can be gained. There is also an increase in supragingival calculus alongside visible nicotine staining. The anterior dentition occasionally have recession and maxillary anterior and palatal surfaces are more adversely affected.
Pathophysiology
Chronic periodontitis is initiated by Gram-negative tooth-associated microbial biofilms that elicit a host response, which results in bone and soft tissue destruction. In response to endotoxin derived from periodontal pathogens, several osteoclast-related mediators target the destruction of alveolar bone and supporting connective tissue such as the periodontal ligament. Major drivers of this aggressive tissue destruction are matrix metalloproteinases (MMPs), cathepsins, and other osteoclast-derived enzymes.
Plaque hypothesis
At least two mechanisms of the microbiology of periodontitis have been described: the specific plaque hypothesis and the non-specific plaque hypothesis.
Consensus is that neither view is entirely correct, but via a middle path, that damage is due to a shift in the relative populations of more and less dangerous bacteria in the plaque. This is called the ecological plaque hypothesis.
The disease is associated with a variable microbial pattern.
Anaerobic species of bacteria Porphyromonas gingivalis, Bacteroides forsythus, Treponema denticola, Prevotella intermedia, Fusobacterium nucleatum, Eubacterium sp. have all been implicated in chronic periodontitis.
Microaerophile bacteria Actinomyces actinomycetemcomitans, Campylobacter rectus, and Eikenella corrodens also may play a role in chronic periodontitis.
Signs and symptoms
In the early stages, chronic periodontitis has few symptoms and in many individuals the disease has progressed significantly before they seek treatment.
Symptoms may include the following:
Redness or bleeding of gums while brushing teeth, using dental floss or biting into hard food (e.g. apples) (though this may occur even in gingivitis, where there is no attachment loss)
Gum swelling that reoccurs
Halitosis, or bad breath, and a persistent metallic taste in the mouth
Gingival recession, resulting in apparent lengthening of teeth. (This may also be caused by heavy-handed brushing or with a stiff tooth brush.)
Deep pockets between the teeth and the gums (pockets are sites where the attachment has been gradually destroyed by collagen-destroying enzymes, known as collagenases)
Loose teeth, in the later stages (though this may occur for other reasons as well)
Drifting of incisors
Gingival inflammation and bone destruction are often painless. Patients sometimes assume that painless bleeding after teeth cleaning is insignificant, although this may be a symptom of progressing chronic periodontitis in that patient.
Subgingival calculus is a frequent finding as well as supragingival calculus due to the bacteria migrating apically and the combined effect of the host response system of the body.
Diagnosis
1999 classification
Chronic periodontitis is one of the seven destructive periodontal diseases as listed in the 1999 classification. Not every case of gingivitis will progress onto chronic periodontitis, but all chronic periodontitis results from gingivitis. Therefore it is important to control the first step; gingival inflammation.
A difficulty which arises with diagnosing chronic periodontitis patients is due to the slow and painless progression of the disease. The most effective and timely diagnosis would be during the mild to moderate stage. However, usually when presenting complaints do arise the effects of mobility and alveolar bone loss have become severe.
A full mouth examination and recording is required to document and track periodontal disease including:
Pocket Depth (PD)
Clinical Attachment Loss (CAL)
Bleeding On Probing (BOP)
Plaque index/score
Furcation involvement
Suppuration
Mobility
Radiographs
Measuring disease progression is carried out by measuring probing pocket depth (PPD) and bleeding indices using a periodontal probe. Pockets greater than 3mm in depth are considered to be unhealthy. True pocket formation of 4 mm or more are specifically related to chronic periodontitis. Bleeding on probing is considered to be a sign of active disease. Discharge of pus, involvement of the root furcation area and deeper pocketing may all indicate reduced prognosis for an individual tooth.
Evidence of alveolar bone loss is also required to differentiate between true bone loss and not attributions of gingival oedema. Usually, a horizontal pattern of bone loss would be found however, vertical (infrabony) bone loss may also be present on specific sites. A Basic Periodontal Examination (BPE) or Periodontal Screening and Recording (PSR) should give a score of 3 or 4. A correct diagnosis is vital to allow the formation of a specific treatment plan for the patient and to arrest the disease progression.
Chronic periodontitis can be further classified into:
Extent (can be either localised affecting < 30% of sites; or generalised if > 30% of sites are affected)
Severity (slight = 1–2 mm CAL; moderate = 3–4 mm CAL; severe ≥5 mm CAL)
2017 classification
Chronic periodontitis is not included within the newer 2017 World Workshop classification. The 2017 classification of Periodontal Diseases and Conditions includes:
Periodontitis:
Necrotizing periodontal diseases
Periodontitis
Periodontitis as a manifestation of systemic disease
Therefore, in accordance to the 2017 classification, a diagnosis would be achieved through the patient assessment individually on the basis of:
Type
Distribution: localised (up to 30% of teeth) or generalised (more than 30% of teeth) and the molar/incisor pattern
Stage and grading
Stages: I (early/mild) with <15% or <2mm interproximal bone loss, II (moderate) with coronal third of root bone loss, III (severe) with mid third of root bone loss, IV (very severe) with apical third of root bone loss
Grading: A (slow) with <0.5% bone loss/age, B (moderate) with 0.5-1.0% bone loss/age, C (rapid) with >1.0% bone loss/age
Status: stable, remission or unstable (see Table 1)
Risk factors, which include systemic diseases such as diabetes or extrinsic factors such as smoking.
Treatment
There is professional agreement among dentists that smoking cessation and good oral hygiene are key to effective treatment and positive outcomes for patients. Similarly, any plaque retentive factors which exist and are modifiable should be corrected, such as overhangs on restorations.
Treatment can involve both non-surgical and surgical therapies. The typical initial treatment known to be effective is scaling and root planing (SRP) to mechanically debride the depths of the periodontal pocket and disrupt the biofilm present. This is done using a powered ultrasonic or sonic scaler and/or unpowered hand instruments.
"In patients with chronic periodontitis, subgingival debridement (in conjunction with supragingival plaque control) is an effective treatment in reducing probing pocket depth and improving the clinical attachment level. In fact it is more effective than supragingival plaque control alone". It is important for patients to be reviewed within 8–12 weeks to assess the treatment response.
Full mouth disinfection protocols are favoured by some clinicians. There is no evidence that full mouth disinfection or full mouth scaling protocols improve the outcome when compared to standard mechanical scaling and root planing.
Open flap debridement
Open flap debridement is used by some practitioners particularly in deeper pocket areas. The advantages of this approach is better visualization of the root surface to be cleaned. This must be weighed against the risks of surgery. Open flap surgery is more effective than non-surgical periodontal therapy in deep pocketing : "Both scaling and root planing alone and scaling and root planing combined with flap procedure are effective methods for the treatment of chronic periodontitis in terms of attachment level gain and reduction in gingival inflammation. In the treatment of deep pockets open flap debridement results in greater PPD reduction and clinical attachment gain."
Guided tissue regeneration
Guided tissue regeneration (GTR) using PTFE membranes is favoured by some practitioners, despite its cost and complexity: "GTR has a greater effect on probing measures of periodontal treatment than open flap debridement, including improved attachment gain, reduced pocket depth, less increase in gingival recession and more gain in hard tissue probing at re-entry surgery. However there is marked variability between studies and the clinical relevance of these changes is unknown. As a result, it is difficult to draw general conclusions about the clinical benefit of GTR. Whilst there is evidence that GTR can demonstrate a significant improvement over conventional open flap surgery, the factors affecting outcomes are unclear from the literature and these might include study conduct issues such as bias. Therefore, patients and health professionals need to consider the predictability of the technique compared with other methods of treatment before making final decisions on use."
Enamel matrix derivative
Enamel matrix derivative (EMD) is favoured by some practitioners despite its high cost: "One year after its application, EMD significantly improved probing attachment levels (1.1 mm) and probing pocket depth reduction (0.9 mm) when compared to a placebo or control, however, the high degree of heterogeneity observed among trials suggests that results have to be interpreted with great caution. In addition, a sensitivity analysis indicated that the overall treatment effect might be overestimated. The actual clinical advantages of using EMD are unknown. With the exception of significantly more postoperative complications in the GTR group, there was no evidence of clinically important differences between GTR and EMD. Bone substitutes may be associated with less gingival recession than EMD."
However, studies have shown that regardless of the conventional periodontal treatments, 20-30% of chronic periodontitis patients do not respond favorably to their treatment. There are many factors which account for these including: ineffective removal of calculus, defective restorations, impaired immune response as a result of a systemic condition, poor plaque control, smoking, etc.
Adjunctive systemic antibiotic treatment
Systemic antibiotics such as amoxicillin or metronidazole are sometimes used in addition to debridement based treatments.
"Systemic antimicrobials in conjunction with scaling and root planing (SRP), can offer an additional benefit over SRP alone in the treatment of periodontitis, in terms of clinical attachment loss (CAL) and probing pocket depth (PPD) change, and reduced risk of additional CAL loss. However, differences in study methodology and lack of data precluded an adequate and complete pooling of data for a more comprehensive analyses. It was difficult to establish definitive conclusions, although patients with deep pockets, progressive or 'active' disease, or specific microbiological profile, can benefit more from this adjunctive therapy."
There is currently low-quality evidence suggesting if adjunctive systemic antimicrobials are beneficial for the non-surgical treatment of periodontitis. It is not sure whether some antibiotics are better than others when used alongside scaling and root planning).
Locally delivered adjunctive antimicrobial treatment
Chemical antimicrobials may be used by the clinician to help reduce the bacterial load in the diseased pocket.
"Among the locally administered adjunctive antimicrobials, the most positive results occurred for tetracycline, minocycline, metronidazole, and chlorhexidine. Adjunctive local therapy generally reduced PD levels....Whether such improvements, even if statistically significant, are clinically meaningful remains a question."
Minocycline is typically delivered via slim syringe applicators.
Chlorhexidine impregnated chips are also available.
Hydrogen peroxide is a naturally occurring antimicrobial that can be delivered directly to the gingival sulcus or periodontal pocket using a custom formed medical device called a Perio Tray. [Title = Custom Tray Application of Peroxide Gel as an Adjunct to Scaling and Root Planing in the Treatment of Periodontitis:
A Randomized, Controlled Three-Month Clinical Trial J Clin Dent 2012;23:48–56.]
Hydrogen peroxide gel was demonstrated to be effective in controlling the bacteria biofilm [Subgingival Delivery of Oral Debriding Agents: A Proof of Concept J Clin Dent 2011;22:149–158] The research shows that a direct application of hydrogen peroxide gel killed virtually all of the bacterial biofilm, was directly and mathematically delivered up to 9mm into periodontal pockets.
Modulating the host response
Sub-antimicrobial doses of doxycycline (SDD) have been used to alter host response to the periodontal pathogens. This is believed to disrupt the action of matrix metalloproteinases and thus minimise host mediated tissue destruction.
"The adjunctive use of SDD with SRP is statistically more effective than SRP alone in reducing PD and in achieving CAL gain."
Systemic Factors
Chronic periodontitis is an inflammatory immune response against the presence of bacteria present. Recent research has suggested that epithelial lining ulceration in chronic periodontal pockets are due to systemic bacterial dissemination and widespread bacterial inflammatory markers present in the host. Two of the most widely investigated systemic diseases associated with chronic periodontitis is diabetes mellitus and cardiovascular disease.
Diabetes Mellitus
Both type 1 and type 2 diabetes have shown a link with the treatment and progression of chronic periodontitis. Chronic periodontitis is more severe in patients that have diabetes than those without, confirming a significant association. With type 2 diabetes patients being shown to have 3.8 times more bone loss and 2.8 times more clinical attachment loss than non-diabetic individuals. With patients with poorly controlled diabetes having a higher risk of alveolar bone loss. Chronic periodontitis can also be a metabolic stressor influencing diabetes control, influencing insulin resistance or becoming a source of inflammatory marker secretion which may strengthen the amount of advanced glycation end product (AGE) mediated cytokine response. Monocytic hyperresponsiveness to bacterial antigen is a biological mechanism that links periodontal disease and diabetes. Increased production of proinflammatory cytokines and mediators cause tissue destruction, attachment loss as well as bone loss causing delayed wound healing.
Cardiovascular Disease
Chronic periodontitis is a marker for cardiovascular disease (CVD). Mechanisms associated with cardiovascular risk are that chronic periodontitis increases inflammatory mediator levels and this may contribute to the onset of CVD, while treatment of chronic periodontitis reduces systemic levels of inflammatory mediators. Certain bacteria found in the periodontal pockets have also been associated to cause atheromatous plaques. Treatment protocol for chronic periodontitis with CVD does not need to be modified as normal periodontal treatment techniques are seen to be effective in CVD patients with additional supportive therapy.
Costs of treatment
"Costs for tooth retention via supportive periodontal therapy are relatively low compared with alternatives (e.g. implants or bridgework) even in periodontally impaired teeth.". However, health outcomes of periodontal therapy are not directly comparable with those from implants or bridgework.
Research
Management
For adults without severe periodontitis and who get routine dental care, regular scale and polish treatment does not make any difference to gingivitis, probing depths or other oral health-related problems. It seems that there is also no difference in plaque levels.
Lasers are increasingly being used in treatments for chronic periodontitis. However, there is some controversy over their use:
"No consistent evidence supports the efficacy of laser treatment as an adjunct to non-surgical periodontal treatment in adults with chronic periodontitis."
Tentative links to other conditions
There is little evidence linking progression of periodontal disease to low birth weight or preterm birth:
"In these women with periodontitis and within this study's limitations, disease progression was not associated with an increased risk for delivering a pre-term or a low birthweight infant."
There is recently emerged evidence linking chronic periodontitis with head and neck squamous cell carcinoma: "Patients with periodontitis were more likely to have poorly differentiated oral cavity SCC than those without periodontitis (32.8% versus 11.5%; P = 0.038)".
| Biology and health sciences | Specific diseases | Health |
1903166 | https://en.wikipedia.org/wiki/Loggerhead%20sea%20turtle | Loggerhead sea turtle | The loggerhead sea turtle (Caretta caretta) is a species of oceanic turtle distributed throughout the world. It is a marine reptile, belonging to the family Cheloniidae. The average loggerhead measures around in carapace length when fully grown. The adult loggerhead sea turtle weighs approximately , with the largest specimens weighing in at more than . The skin ranges from yellow to brown in color, and the shell is typically reddish brown. No external differences in sex are seen until the turtle becomes an adult, the most obvious difference being the adult males have thicker tails and shorter plastrons (lower shells) than the females.
The loggerhead sea turtle is found in the Atlantic, Pacific, and Indian Oceans, as well as the Mediterranean Sea. It spends most of its life in saltwater and estuarine habitats, with females briefly coming ashore to lay eggs. The loggerhead sea turtle has a low reproductive rate; females lay an average of four egg clutches and then become quiescent, producing no eggs for two to three years. The loggerhead reaches sexual maturity within 17–33 years and has a lifespan of 47–67 years.
The loggerhead sea turtle is omnivorous, feeding mainly on bottom-dwelling invertebrates. Its large and powerful jaws serve as an effective tool for dismantling its prey. Young loggerheads are exploited by numerous predators; the eggs are especially vulnerable to terrestrial organisms. Once the turtles reach adulthood, their formidable size limits predation to large marine animals, such as large sharks.
The loggerhead sea turtle is considered a vulnerable species by the International Union for Conservation of Nature.
In total, nine distinct population segments are under the protection of the Endangered Species Act of 1973, with four population segments classified as "threatened" and five classified as "endangered".
Commercial international trade of loggerheads or derived products is prohibited by CITES Appendix I.
Untended fishing gear is responsible for many loggerhead deaths. The greatest threat is loss of nesting habitat due to coastal development, predation of nests, and human disturbances (such as coastal lighting and housing developments) that cause disorientations during the emergence of hatchlings. Turtles may also suffocate if they are trapped in fishing trawls. Turtle excluder devices have been implemented in efforts to reduce mortality by providing an escape route for the turtles. Loss of suitable nesting beaches and the introduction of exotic predators have also taken a toll on loggerhead populations. Efforts to restore their numbers will require international cooperation, since the turtles roam vast areas of ocean and critical nesting beaches are scattered across several countries.
Taxonomy
Carl Linnaeus gave the loggerhead its first binomial name, Testudo caretta, in 1758. Thirty-five other names emerged over the following two centuries, with the combination Caretta caretta first introduced in 1873 by Leonhard Stejneger. The English common name "loggerhead" refers to the animal's large head. The loggerhead sea turtle belongs to the family Cheloniidae, which includes all extant sea turtles except the leatherback sea turtle. The subspecific classification of the loggerhead sea turtle is debated, but most authors consider it a single polymorphic species. Molecular genetics has confirmed hybridization of the loggerhead sea turtle with the Kemp's ridley sea turtle, hawksbill sea turtle, and green sea turtles. The extent of natural hybridization is not yet determined; however, second-generation hybrids have been reported, suggesting some hybrids are fertile.
Evolution
Although evidence is lacking, modern sea turtles probably descended from a single common ancestor during the Cretaceous period. Like all other sea turtles except the leatherback, loggerheads are members of the ancient family Cheloniidae, and appeared about 40 million years ago. Of the six species of living Cheloniidae, loggerheads are more closely related to the Kemp's ridley sea turtle, olive ridley sea turtle, and the hawksbill turtle than they are to the flatback turtle and the green turtle.
Around three million years ago, during the Pliocene epoch, Central America emerged from the sea, effectively cutting off currents between the Atlantic and Indo-Pacific Oceans. The rerouting of ocean currents led to climatic changes as the Earth entered a glacial cycle. Cold water upwelling around the Cape of Good Hope and reduction in water temperature at Cape Horn formed coldwater barriers to migrating turtles. The result was a complete isolation of the Atlantic and Pacific populations of loggerheads. During the most recent ice age, the beaches of southeastern North America were too cold for sea turtle eggs. As the Earth began to warm, loggerheads moved farther north, colonizing the northern beaches. Because of this, turtles nesting between North Carolina and northern Florida represent a different genetic population from those in southern Florida.
The distinct populations of loggerheads have unique characteristics and genetic differences. For example, Mediterranean loggerheads are smaller, on average, than Atlantic Ocean loggerheads. North Atlantic and Mediterranean loggerhead sea turtles are descendants of colonizing loggerheads from Tongaland, South Africa. South African loggerhead genes are still present in these populations today.
Description
The loggerhead sea turtle is the world's largest hard-shelled turtle, slightly larger at average and maximum mature weights than the green sea turtle and the Galapagos tortoise. It is also the world's second largest extant turtle after the leatherback sea turtle. Adults have an approximate weight range of , averaging around , and a straight-line carapace length range of . The maximum reported weight is and the maximum (presumed total) length is . The head and carapace (upper shell) range from a yellow-orange to a reddish brown, while the plastron (underside) is typically pale yellow. The turtle's neck and sides are brown on the tops and yellow on the sides and bottom.
The turtle's shell is divided into two sections: carapace and plastron. The carapace is further divided into large plates, or scutes. Typically, 11 or 12 pairs of marginal scutes rim the carapace. Five vertebral scutes run down the carapace's midline, while five pairs of costal scutes border them. The nuchal scute is located at the base of the head. The carapace connects to the plastron by three pairs of inframarginal scutes forming the bridge of the shell. The plastron features paired gular, humeral, pectoral, abdominal, femoral, and anal scutes. The shell serves as external armor, although loggerhead sea turtles cannot retract their heads or flippers into their shells.
Sexual dimorphism of the loggerhead sea turtle is only apparent in adults. Adult males have longer tails and claws than females. The males' plastrons are shorter than the females', presumably to accommodate the males' larger tails. The carapaces of males are wider and less domed than the females', and males typically have wider heads than females. The sex of juveniles and subadults cannot be determined through external anatomy, but can be observed through dissection, laparoscopy (an operation performed on the abdomen), histological examination (cell anatomy), and radioimmunological assays (immune study dealing with radiolabeling).
Lachrymal glands located behind each eye allow the loggerhead to maintain osmotic balance by eliminating the excess salt obtained from ingesting ocean water. On land, the excretion of excess salt gives the false impression that the turtle is crying. The urea content is high in Caretta caretta tears.
The skull is most easily distinguished from other sea turtles by having maxillae that meet in the mid-line of the palate. The portion of skull behind the eyes is also relatively large and bulbous due to the extensive jaw muscles.
Distribution
The loggerhead sea turtle has a cosmopolitan distribution, nesting over the broadest geographical range of any sea turtle. It inhabits the Atlantic, Indian, and Pacific Oceans and the Mediterranean Sea.
In the Atlantic Ocean, the greatest concentration of loggerheads is along the southeastern coast of North America and in the Gulf of Mexico. Very few loggerheads are found along the European and African coastlines. Florida is the most popular nesting site, with more than 67,000 nests built per year. Nesting extends as far north as Virginia, as far south as Brazil, and as far east as the Cape Verde Islands. The Cape Verde Islands are the only significant nesting site on the eastern side of the Atlantic. Loggerheads found in the Atlantic Ocean feed from Canada to Brazil.
In the Indian Ocean, loggerheads feed along the coastlines of Africa, the Arabian Peninsula, and in the Arabian Sea. Along the African coastline, loggerheads nest from Mozambique's Bazaruto Archipelago to South Africa's St Lucia estuary. The largest Indian Ocean nesting site is Oman, on the Arabian Peninsula, which hosts around 15,000 nests, giving it the second largest nesting population of loggerheads in the world. Western Australia is another notable nesting area, with 1,000–2,000 nests per year.
Pacific loggerheads live in temperate to tropical regions. They forage in the East China Sea, the southwestern Pacific, and along the Baja California Peninsula. Eastern Australia and Japan are the major nesting areas, with the Great Barrier Reef deemed an important nesting area. Pacific loggerheads occasionally nest in Vanuatu and Tokelau. Yakushima Island is the most important site, with three nesting grounds visited by 40% of all nearby loggerheads. After nesting, females often find homes in the East China Sea, while the Kuroshio Current Extension's Bifurcation region provides important juvenile foraging areas. Eastern Pacific populations are concentrated off the coast of Baja California, where upwelling provides rich feeding grounds for juvenile turtles and subadults. Nesting sites along the eastern Pacific Basin are rare. mtDNA sequence polymorphism analysis and tracking studies suggest 95% of the population along the coast of the Americas hatch on the Japanese Islands in the western Pacific. The turtles are transported by the prevailing currents across the full length of the northern Pacific, one of the longest migration routes of any marine animal. The return journey to the natal beaches in Japan has been long suspected, although the trip would cross unproductive clear water with few feeding opportunities. Evidence of a return journey came from an adult female loggerhead named Adelita, which in 1996, equipped with a satellite tracking device, made the trip from Mexico across the Pacific. Adelita was the first animal of any kind ever tracked across an ocean basin.
The Mediterranean Sea is a nursery for juveniles, as well as a common place for adults in the spring and summer months. Almost 45% of the Mediterranean juvenile population has migrated from the Atlantic. Loggerheads feed in the Alboran Sea and the Adriatic Sea, with tens of thousands of specimens (mainly sub-adult) seasonally present in the North-Eastern portion of the latter, above all in the area of the Po Delta.
Greece is the most popular nesting site along the Mediterranean, with more than 3,000 nests per year. Zakynthos hosts the largest Mediterranean nesting with the second one being in Kyparissia Bay. Because of this, Greek authorities do not allow planes to take off or land at night in Zakynthos due to the nesting turtles. In addition to the Greek coast, the coastlines of Cyprus and Turkey are also common nesting sites.
One record of this turtle was made in Ireland when a specimen washed ashore on Ballyhealy Beach in County Wexford. Another records one specimen being washed up on a beach in County Donegal, Ireland.
Habitat
Loggerhead sea turtles spend most of their lives in the open ocean and in shallow coastal waters. They rarely come ashore besides the females' brief visits to construct nests and deposit eggs. Hatchling loggerhead turtles live in floating mats of Sargassum algae. Adults and juveniles live along the continental shelf as well as in shallow coastal estuaries. In the northwestern Atlantic Ocean, age plays a factor in habitat preference. Juveniles are more frequently found in shallow estuarine habitats with limited ocean access compared to non-nesting adults. Loggerheads occupy waters with surface temperatures ranging from during non-nesting season. Temperatures from are most suitable for nesting females.
Juvenile loggerheads share the Sargassum habitat with a variety of other organisms. The mats of Sargassum contain as many as 100 different species of animals on which the juveniles feed. Prey found in Sargassum mats may include barnacles, crab larvae, fish eggs, and hydrozoan colonies. Some prey, such as ants, flies, aphids, leafhoppers, and beetles, are carried by the wind to the mats. Marine mammals and commercial fishes, including tuna and mahi-mahi, also inhabit the Sargassum mats.
Behavior
Loggerhead sea turtles observed in captivity and in the wild are most active during the day. In captivity, the loggerheads' daily activities are divided between swimming and resting on the bottom. While resting, they spread their forelimbs to about midstroke swimming position. They remain motionless with eyes open or half-shut and are easily alerted during this state. At night, captives sleep in the same position with their eyes tightly shut, and are slow to react. Loggerheads spend up to 85% of their day submerged, with males being the more active divers than females. The average duration of dives is 15–30 min, but they can stay submerged for up to four hours. Juvenile loggerheads and adults differ in their swimming methods. A juvenile keeps its forelimbs pressed to the side of its carapace, and propels itself by kicking with its hind limbs. As the juvenile matures, its swimming method is progressively replaced with the adult's alternating-limb method. They depend entirely on this method of swimming by one year old.
Water temperature affects the sea turtle's metabolic rate. Lethargy is induced at temperatures between . The loggerhead takes on a floating, cold-stunned posture when temperatures drop to around . However, younger loggerheads are more resistant to cold and do not become stunned until temperatures drop below . The loggerheads' migration helps to prevent instances of cold-stunning. Higher water temperatures cause an increase in metabolism and heart rate. A loggerhead's body temperature increases in warmer waters more quickly than it decreases in colder water; their critical thermal maximum is currently unknown. In February 2015, a live loggerhead turtle was found floating in British Columbian waters of with extensive algal growth on its carapace.
Female-female aggression, which is fairly rare in other marine vertebrates, is common among loggerheads. Ritualized aggression escalates from passive threat displays to combat. This conflict primarily occurs over access to feeding grounds. Escalation typically follows four steps. First, initial contact is stimulated by visual or tactile cues. Second, confrontation occurs, beginning with passive confrontations characterized by wide head-tail circling. They begin aggressive confrontation when one turtle ceases to circle and directly faces the other. Third, sparring occurs with turtles snapping at each other's jaws. The final stage, separation, is either mutual, with both turtles swimming away in opposite directions, or involves chasing one out of the immediate vicinity. Escalation is determined by several factors, including hormone levels, energy expenditure, expected outcome, and importance of location. At all stages, an upright tail shows willingness to escalate, while a curled tail shows willingness to submit. Because higher aggression is metabolically costly and potentially debilitating, contact is much more likely to escalate when the conflict is over access to good foraging grounds. Further aggression has also been reported in captive loggerheads. The turtles are seemingly territorial, and will fight with other loggerheads and sea turtles of different species.
Feeding
The loggerhead sea turtle is omnivorous, feeding mainly on bottom-dwelling invertebrates, such as gastropods, bivalves, decapods, and horseshoe crabs. It has a greater list of known prey than any other sea turtle. Other food items include sponges, corals, sea pens, polychaete worms, tube worms, sea anemones, cephalopods, barnacles, brachiopods, amphipods, isopods, Portuguese men o' war, insects, bryozoans, hydrozoans, sea urchins, sand dollars, sea cucumbers, starfish, tunicates, fish (eggs, juveniles, and adults), hatchling turtles (including members of its own species), algae, and vascular plants. During migration through the open sea, loggerheads eat jellyfish, floating molluscs, floating egg clusters, squid, and flying fish.
Loggerheads crush prey with their large and powerful jaws. Projecting scale points on the anterior margin of the forelimbs allow manipulation of the food. These points can be used as "pseudo-claws" to tear large pieces of food in the loggerhead's mouth. The loggerhead will turn its neck sideways to consume the torn food on the scale points. Inward-pointing, mucus-covered papillae found in the fore region of the loggerhead's esophagus filter out foreign bodies, such as fish hooks. The next region of the esophagus is not papillated, with numerous mucosal folds. The digestion rate in loggerheads is temperature-dependent; it increases as temperature increases.
Predators
Loggerheads have numerous predators, especially early in their lives. Egg and nestling predators include ghost crabs, oligochaete worms, some beetles, flesh fly larvae, some ants, flesh flies, snakes, gulls, corvids, opossums, bears, rats, armadillos, mustelids, skunks, canids like coyotes, dingos, the Red foxes in Australia, Jackals and feral dogs, procyonids, Feral cats, Feral pigs, and humans. During their migration from their nests to the sea, hatchlings are preyed on by dipteran larvae, crabs, toads, lizards, snakes, seabirds such as frigatebirds, and other assorted birds and mammals. In the ocean, predators of the loggerhead juveniles include portunid crabs and various fishes, such as parrotfishes and moray eels. Adults are more rarely attacked due to their large size, but may be preyed on by large sharks, seals, and killer whales. Nesting females are attacked by flesh flies, feral dogs, and humans. Salt marsh mosquitos can also pester nesting females.
In Australia, the introduction of the red fox (Vulpes vulpes) by British settlers in the 19th century led to significant reductions in loggerhead sea turtle populations. In one coastal section in eastern Australia during the 1970s, predation of turtle eggs destroyed up to 95% of all clutches laid. Aggressive efforts to destroy foxes in the 1980s and 1990s has reduced this impact; however, it is estimated that it will be the year 2020 before populations will experience complete recovery from such dramatic losses.
Along the southeastern coast of the United States, the raccoon (Procyon lotor) is the most destructive predator of nesting sites. Mortality rates of nearly 100% of all clutches laid in a season have been recorded on some Florida beaches. This is attributed to an increase in raccoon populations, which have flourished in urban environments. Aggressive efforts to protect nesting sites by covering them with wire mesh has significantly reduced the impact of raccoon predation on loggerhead sea turtle eggs.
Up to 40% of nesting females around the world have wounds believed to come from shark attacks.
Disease and parasites
Infectious bacteria such as Pseudomonas and Salmonella attack loggerhead hatchlings and eggs. Fungi such as Penicillium infect loggerhead sea turtle nests and cloacae.
Fibropapillomatosis disease caused by a form of the herpes-type virus threatens loggerheads with internal and external tumors. These tumors disrupt essential behaviors and, if on the eyes, cause permanent blindness. Trematodes of the family Spirorchiidae inhabit tissues throughout the body of the loggerhead, including vital organs, such as the heart and the brain. Trematode infection can be highly debilitating. For example, inflammatory trematode lesions can cause endocarditis and neurological disease. A nematode, Angiostoma carettae, also infects loggerheads, causing histologic lesions in the respiratory tract.
More than 100 species of animals from 13 phyla, as well as 37 kinds of algae, live on loggerheads' backs. These parasitic organisms, which increase drag, offer no known benefit to the turtle, although the dulling effect of organisms on shell color may improve camouflage.
In 2018, researchers from Florida State University examined 24 individual turtle carapaces and found an average of 33,000 meiofauna with one turtle having 150,000 organisms living on the shell. A collection of 7,000 nematodes from 111 genera were found on the turtles studied.
Life history
Early life
Hatchlings range in color from light brown to almost black, lacking the adult's distinct yellows and reds. Upon hatching, they measure about and weigh about . The eggs are typically laid on the beach in an area above the high-tide line. The eggs are laid near the water so the hatchlings can return to the sea. The loggerhead's sex is dictated by the temperature of the underground nest. Incubation temperatures generally range from . Sea turtle eggs kept at a constant incubating temperature of 32 °C become females. Eggs incubating at 28 °C become males. An incubation temperature of 30 °C results in an equal ratio of male to female hatchlings. Hatchlings from eggs in the middle of the clutch tend to be the largest, grow the fastest, and be the most active during the first few days of sea life.
After incubating for around 80 days, hatchlings dig through the sand to the surface, usually at night, when darkness increases the chance of escaping predation and damage from extreme sand surface temperatures is reduced. Hatchlings enter the ocean by navigating toward the brighter horizon created by the reflection of the moon and starlight off the water's surface.
Hatchlings can lose up to 20% of their body mass due to evaporation of water as they journey from nest to ocean. They initially use the undertow to push them five to 10 m away from the shore. Once in the ocean, they swim for about 20 hours, taking them far offshore. An iron compound, magnetite, in their brains allows the turtles to perceive the Earth's magnetic field, for navigation.
Many hatchlings use Sargassum in the open ocean as protection until they reach . Hatchling loggerheads live in this pelagic environment until they reach juvenile age, and then they migrate to nearshore waters.
Maturation
When ocean waters cool, loggerheads must migrate to warmer areas or hibernate to some degree. In the coldest months, they submerge for up to seven hours at a time, emerging for only seven minutes to breathe. Although outdone by freshwater turtles, these are among the longest recorded dives for any air-breathing marine vertebrate. During their seasonal migration, juvenile loggerheads have the ability to use both magnetic and visual cues. When both aids are available, they are used in conjunction; if one aid is not available, the other suffices. The turtles swim at about during migration.
Like all marine turtles, the loggerhead prepares for reproduction in its foraging area. This takes place several years before the loggerhead migrates to a mating area. Female loggerheads first reproduce at ages 28–33 in Southeastern United States and Australia, and at ages 17–30 in South Africa. Age at first reproduction in the Mediterranean, Oman, Japan, and Brazil are unknown. Nesting loggerheads have a straight carapace length of . Because of the large range, carapace length is not a reliable indicator of sexual maturity. Their estimated maximum lifespan is 47–67 years in the wild.
Reproduction
Female loggerheads first reproduce between the ages of 17 and 33, and their mating period may last more than six weeks. They court their mates, but these behaviors have not been thoroughly examined. Male forms of courtship behavior include nuzzling, biting, and head and flipper movements. Studies suggest females produce cloacal pheromones to indicate reproductive ability. Before mating, the male approaches a female and attempts to mount her, while she resists. Next, the male and female begin to circle each other. If the male has competitors, the female may let the males struggle with each other. The winner then mounts the female; the male's curved claws usually damage the shoulders of the female's shell during this process. Other courting males bite the male while he is attempting to copulate, damaging his flippers and tail, possibly exposing bones. Such damage can cause the male to dismount and may require weeks to heal.
While nesting, females produce an average of 3.9 egg clutches, and then become quiescent, producing no eggs for two to three years. Unlike other sea turtles, courtship and mating usually do not take place near the nesting beach, but rather along migration routes between feeding and breeding grounds. Recent evidence indicates ovulation in loggerheads is mating-induced. Through the act of mating, the female ovulates eggs which are fertilized by the male. This is unique, as mating-induced ovulation is rare outside of mammals. In the Northern Hemisphere, loggerheads mate from late March to early June. The nesting season is short, between May and August in the Northern Hemisphere and between October and March in the Southern Hemisphere.
Loggerheads may display multiple paternity. Multiple paternity is possible due to sperm storage. The female can store sperm from multiple males in her oviducts until ovulation. A single clutch may have as many as seven fathers, each contributing sperm to a portion of the clutch. Multiple paternity and female size are positively correlated. Two hypotheses explain this correlation. One posits that males favor large females because of their perceived higher fecundity (ability to reproduce). The other states, because larger females are able to swim more quickly to mating grounds, they have longer mating periods.
All sea turtles have similar basic nesting behaviors. Females return to lay eggs at intervals of 12–17 days during the nesting season, on or near the beach where they hatched. They exit the water, climb the beach, and scrape away the surface sand to form a body pit. With their hind limbs, they excavate an egg chamber in which the eggs are deposited. The females then cover the egg chamber and body pit with sand, and finally return to the sea. This process takes one to two hours, and occurs in open sand areas or on top of sand dunes, preferably near dune grasses that the females can use to camouflage the nest. The nesting area must be selected carefully because it affects characteristics such as fitness, emergence ratio, and vulnerability to nest predators. Loggerheads have an average clutch size of 112.4 eggs.
Conservation
Many human activities have negative effects on loggerhead sea turtle populations. The prolonged time required for loggerheads to reach sexual maturity and the high mortality rates of eggs and young turtles from natural phenomena compound the problems of population reduction as a consequence of human activities.
Threats
Loggerhead sea turtles were once intensively hunted for their meat and eggs; consumption has decreased, however, due to worldwide legislation. Despite this, turtle meat and eggs are still consumed in countries where regulations are not strictly enforced. In Mexico, turtle eggs are a common meal; locals claim the egg is an aphrodisiac. Eating turtle eggs or meat can cause serious illness due to harmful bacteria, such as Pseudomonas aeruginosa and Serratia marcescens, and high levels of toxic metals that build up through bioaccumulation.
The US West Coast is a critical migratory corridor for the Pacific loggerheads, in which these turtles swim across the Pacific to California's coast from breeding grounds in Japan. Important foraging habitats for juveniles in the central North Pacific have been revealed through telemetry studies. Along with these foraging habitats, high levels of bycatch from industrial-scale fisheries have been found to overlap; with drift gillnets in the past and longline fisheries presently. Many juvenile loggerheads aggregate off the coast of Baja California Sur, Mexico, where small coastal fisheries increase these turtles' mortality risk; fishers have reported catching dozens of loggerheads with bottom-set gear per day per boat. The most common commercial fishery that accidentally takes loggerheads are bottom trawls used for shrimp vessels in the Gulf of California. In 2000, between 2,600 and 6,000 loggerheads were estimated to have been killed by pelagic longlining in the Pacific.
Fishing gear is the biggest threat to loggerheads in the open ocean. They often become entangled in longlines or gillnets. According to the 2009 status review of loggerheads by the Fisheries Service, drowning from entanglement in longline and gillnet fishing gear is the turtles' primary threat in the North Pacific. They also become stuck in traps, pots, trawls, and dredges. Caught in this unattended equipment, loggerheads risk serious injury or drowning. Turtle excluder devices for nets and other traps reduce the number being accidentally caught.
Nearly 11 million metric tons of plastic are released into the ocean annually. A number that is projected to increase to 29 million metric tons by 2040. Turtles ingest a wide array of this floating debris, including bags, sheets, pellets, balloons and abandoned fishing line. Loggerheads may mistake the floating plastic for jellyfish, a common food item. The ingested plastic causes numerous health concerns, including intestinal blockage, reduced nutrient absorption and malnutrition, suffocation, ulcerations, or starvation. Ingested plastics release toxic compounds, including polychlorinated biphenyls, which may accumulate in internal tissues. Such toxins may lead to a thinning of eggshells, tissue damage, or deviation from natural behaviors.
Artificial lighting discourages nesting and interferes with the hatchlings' ability to navigate to the water's edge. Females prefer nesting on beaches free of artificial lighting. On developed beaches, nests are often clustered around tall buildings, perhaps because they block out the man-made light sources. Loggerhead hatchlings are drawn toward the brighter area over the water which is the consequence of the reflection of moon and star light. Confused by the brighter artificial light, they navigate inland, away from the protective waters, which exposes them to dehydration and predation as the sun rises. Artificial lighting causes tens of thousands of hatchling deaths per year.
Destruction and encroachment of habitat by humans is another threat to loggerhead sea turtles. Optimum nesting beaches are open-sand beaches above the high-tide line. However, beach development deprives them of suitable nesting areas, forcing them to nest closer to the surf. Urbanization often leads to the siltation of sandy beaches, decreasing their viability. Construction of docks and marinas can destroy near-shore habitats. Boat traffic and dredging degrades habitat and can also injure or kill turtles when boats collide with turtles at or near the surface.
Annual variations in climatic temperatures can affect sex ratios, since loggerheads have temperature-dependent sex determination. High sand temperatures may skew gender ratios in favor of females. Nesting sites exposed to unseasonably warm temperatures over a three-year period produced 87–99% females. This raises concern over the connection between rapid global temperature changes and the possibility of population extinction. A more localized effect on gender skewing comes from the construction of tall buildings, which reduce sun exposure, lowering the average sand temperature, which results in a shift in gender ratios to favor the emergence of male turtles. Construction of new thermal power stations can raise local water temperature, which is also said to be a threat.
The increase of temperature and food availability will increase reproduction output of loggerhead turtles. Many researchers agree that temperature increases due to climate change has a complicated impact on turtles. At breeding sites when a loggerhead turtle lays multiple clutches in a season, a higher temperature will cause the duration of time between laying two different nests to become shorter. The amount of food availability makes a difference in reproductive output because when there is a greater amount of food available, the turtles will grow to a larger size. The larger a turtle is, the more likely they will have a greater reproductive output. The amount of food also has a relationship to temperature. Researchers have found that an increase of temperature causes feeding grounds to produce more food.
Tropical Cyclones have a significant impact on hatchling loss. The associated storm surges push water higher up the beach, flooding nest and drowning the embryos. Strong wave action may eroded away sand, exposing the eggs to drying and predation. The current trend of rising sea surface temperatures and the increase in both numbers and intensities of tropical cyclones as a result of climate change pose a growing threat to turtle populations.
Conservation efforts
Since the loggerhead occupies such a broad range, successful conservation requires efforts from multiple countries.
Loggerhead sea turtles are classified as vulnerable by the International Union for Conservation of Nature and are listed under Appendix I of the Convention on International Trade in Endangered Species, making commercial international trade prohibited. In the United States, the Fish and Wildlife Service and National Marine Fisheries Service classify them as a threatened species under the Endangered Species Act. Loggerheads are listed as endangered under both Australia's Environment Protection and Biodiversity Conservation Act 1999 and Queensland's Nature Conservation Act 1992. The Convention on Migratory Species works for the conservation of loggerhead sea turtles on the Atlantic coast of Africa, as well as in the Indian Ocean and southeast Asia. Throughout Japan, the Sea Turtle Association of Japan aids in the conservation of loggerhead sea turtles. Greece's ARCHELON works for their conservation. The Marine Research Foundation works for loggerhead conservation in Oman. Annex 2 of the Specially Protected Areas and Wildlife Protocol of the Cartagena Convention, which deals with pollution that could harm marine ecosystems, also protects them. Conservation organizations worldwide have worked with the shrimp trawling industry to develop turtle exclusion devices (TEDs) to exclude even the largest turtles. TEDs are mandatory for all shrimp trawlers.
In many places during the nesting season, workers and volunteers search the coastline for nests, and researchers may also go out during the evening to look for nesting females for tagging studies and gather barnacles and tissues samples. Volunteers may, if necessary, relocate the nests for protection from threats, such as high spring tides and predators, and monitor the nests daily for disturbances. After the eggs hatch, volunteers uncover and tally hatched eggs, undeveloped eggs, and dead hatchlings. Any remaining live hatchlings are released or taken to research facilities. Typically, those that lack the vitality to hatch and climb to the surface die.
United States
The National Marine Fisheries Service (NMFS), National Oceanic and Atmospheric Administration (NOAA), the U.S. Fish and Wildlife Services (USFWS), and the Department of the Interior ruled four distinct population segments as threatened (Northwestern Atlantic Ocean, South Atlantic Ocean, Southeast Indo-Pacific Ocean, and Southwest Indian Ocean) and five as endangered (Mediterranean Sea, North Indian Ocean, North Pacific Ocean, Northeast Atlantic Ocean, and South Pacific Ocean) effective on October 24, 2011.
Off the coast of southern California NMFS, NOAA, and Department of Commerce prohibited fishing with large drift gillnet (DGN) gear in the loggerhead conservation area during the presence of El Niño conditions in order to protect the endangered North Pacific Ocean loggerhead DPS. This ruling effective July 23, 2014 was intended to prevent bycatch of loggerhead sea turtles. A team including sea turtle biologists and oceanographers determined the presence of El Niño conditions based on the El Niño watch issued by the Climate Prediction Center (CPC), anomalies found in sea surface temperature (SST) charts published by NOAA's Coast Watch Program, the presence of loggerhead sea turtles in the Pacific loggerhead conservation area, and reports of loggerhead strandings. The SST data showed higher than average temperatures during summer months off the coast of southern California. This same fisheries closure ruling due to El Niño conditions was again implemented May 29, 2015, and then again June 1, 2016.
Critical habitat designation for the Northwest Atlantic Ocean DPS of loggerhead sea turtles specified 38 marine areas that include nearshore reproductive habitat, breeding areas, winter area, constricted migratory corridors, and Sargassum habitat. This ruling was made the NMFS, NOAA, and Department of Commerce effective August 11, 2014. Nesting beaches were identified as critical terrestrial habitat by Fish and Wildlife Services and the Department of the Interior within the Atlantic Ocean and Gulf of Mexico, effective August 11, 2014. The 2012 BiOp is an integral component to managing the shallow-set fishery, because the one-year incidental take statement (ITS, including reasonable and prudent management measures, and terms and conditions) forms the basis for regulations that specify the annual limits on leatherback and North Pacific loggerhead sea turtle interactions with the fishery that are necessary to manage the impacts of the fishery on sea turtles.
Effective January 11, 2010 the NMFS, NOAA, and Department of Commerce removed the limit on the number of fishing gear deployments for the Hawaii-based pelagic shallow-set longline fisheries and simultaneously increased the number of incidental interactions allowed with loggerhead sea turtles. This ruling stated that longline fisheries may not interact with over 46 loggerhead sea turtles a year, a number thought to not interfere with survival and recovery of loggerhead sea turtles. This ruling was revised March 10, 2011 to reduce the number of allowed interactions from 46 a year to 17, a revision aimed to protect the loggerheads and maintain fishery yield. November 18, 2011 the pelagic shallow-set longline fisheries in Hawaii reached the annual limit on physical interactions with turtles and was closed by NMFS. Incidental interaction limit for loggerhead turtles was increased from 17 to 34 interactions a year starting November 5, 2012.
Symbols
The loggerhead sea turtle appears on the $1000 Colombian peso coin. In the United States, the loggerhead sea turtle is the official state reptile of South Carolina and also the state saltwater reptile of Florida.
| Biology and health sciences | Turtles | Animals |
1904332 | https://en.wikipedia.org/wiki/Brassica%20rapa | Brassica rapa | Brassica rapa is a plant species that has been widely cultivated into many forms, including the turnip (a root vegetable), komatsuna, napa cabbage, bomdong, bok choy, and rapini.
Brassica rapa subsp. oleifera is an oilseed commonly known as turnip rape, field mustard, bird's rape, and keblock. Rapeseed oil is a general term for oil from some Brassica species. Food grade oil made from the seed of low-erucic acid Canadian-developed strains is also called canola oil, while non-food oil is called colza oil. Canola oil can be sourced from Brassica rapa and Brassica napus, which are commonly grown in Canada, and Brassica juncea, which is less common.
History
The geographic and genetic origins of B. rapa have been difficult to identify due to its long history of human cultivation. It is found in most parts of the world, and has returned to the wild many times as a feral plant or weed.
Genetic sequencing and environmental modelling have indicated that ancestral B. rapa likely originated 4000 to 6000 years ago in the Hindu Kush area of Central Asia, and had three sets of chromosomes, providing the genetic potential for a diversity of form, flavour, and growth. Domestication has produced modern vegetables and oil-seed crops, all with two sets of chromosomes.
Oilseed subspecies (subsp. oleifera) of Brassica rapa may have been domesticated several times from the Mediterranean to India, starting as early as 2000 BC. There are descriptions of B. rapa vegetables in Indian and Chinese documents from around 1000 BC.
Edible turnips were possibly first cultivated in northern Europe, and were an important food in ancient Rome. The turnip then spread east to China, and reached Japan by 700 AD.
In the 18th century, the turnip and the oilseed-producing variants were thought to be different species by Carl Linnaeus, who named them B. rapa and B. campestris. Twentieth-century taxonomists found that the plants were cross fertile and thus belonged to the same species. Since the turnip had been named first by Linnaeus, the name Brassica rapa was adopted.
Uses
Many butterflies, including the small white, feed from and pollinate the B. rapa flowers.
The young leaves are a common leaf vegetable and can be eaten raw; older leaves are typically cooked. The taproot and seeds can also be eaten raw, although the seeds contain an oil that can cause irritation for some people.
Cultivars
| Biology and health sciences | Brassicales | Plants |
1904373 | https://en.wikipedia.org/wiki/Transformation%20%28function%29 | Transformation (function) | In mathematics, a transformation, transform, or self-map is a function f, usually with some geometrical underpinning, that maps a set X to itself, i.e. .
Examples include linear transformations of vector spaces and geometric transformations, which include projective transformations, affine transformations, and specific affine transformations, such as rotations, reflections and translations.
Partial transformations
While it is common to use the term transformation for any function of a set into itself (especially in terms like "transformation semigroup" and similar), there exists an alternative form of terminological convention in which the term "transformation" is reserved only for bijections. When such a narrow notion of transformation is generalized to partial functions, then a partial transformation is a function f: A → B, where both A and B are subsets of some set X.
Algebraic structures
The set of all transformations on a given base set, together with function composition, forms a regular semigroup.
Combinatorics
For a finite set of cardinality n, there are nn transformations and (n+1)n partial transformations.
| Mathematics | Geometry: General | null |
1906103 | https://en.wikipedia.org/wiki/Human%20reproductive%20system | Human reproductive system | The human reproductive system includes the male reproductive system, which functions to produce and deposit sperm, and the female reproductive system, which functions to produce egg cells and to protect and nourish the fetus until birth. Humans have a high level of sexual differentiation. In addition to differences in nearly every reproductive organ, there are numerous differences in typical secondary sex characteristics.
Human reproduction usually involves internal fertilization by sexual intercourse. In this process, the male inserts his penis into the female's vagina and ejaculates semen, which contains sperm. A small proportion of the sperm pass through the cervix into the uterus and then into the fallopian tubes for fertilization of the ovum. Only one sperm is required to fertilize the ovum. Upon successful fertilization, the fertilized ovum, or zygote, travels out of the fallopian tube and into the uterus, where it implants in the uterine wall. This marks the beginning of gestation, better known as pregnancy, which continues for around nine months as the fetus develops. When the fetus has developed to a certain point, pregnancy is concluded with childbirth, involving labor. During labor, the uterine muscles contract, and the cervix dilates typically over a period of hours, allowing the infant to pass from the uterus through the vagina. Human infants are entirely dependent on their caregivers and require parental care. Infants rely on their caregivers for comfort, cleanliness, and food. Food may be provided by breastfeeding or formula feeding.
Structure
Female
The human female reproductive system is a series of organs primarily located inside the body and around the pelvic region of a female that contribute towards the reproductive process. The human female reproductive system contains three main parts: the vagina, which leads from the vulva, the vaginal opening, to the uterus; the uterus, which holds the developing fetus; and the ovaries, which produce the female's ova. The breasts are involved during the parenting stage of reproduction, but in most classifications they are not considered to be part of the female reproductive system.
The vagina meets the outside at the vulva, which is made up of the labia, clitoris and vestibule; during intercourse this area is lubricated by mucus secreted by the Bartholin's glands. The vagina is attached to the uterus through the cervix, while the uterus is attached to the ovaries via the fallopian tubes. Each ovary contains hundreds of egg cells or ova (singular ovum).
Approximately every 28 days, the pituitary gland releases a hormone that stimulates some of the ova to develop and grow. One ovum is released and it passes through the fallopian tube into the uterus. Hormones produced by the ovaries prepare the uterus to receive the ovum. The lining of the uterus, called the endometrium, and unfertilized ova are shed each cycle through the process of menstruation. If the ova is fertilized by sperm, it attaches to the endometrium and the fetus develops.
Male
The male reproductive system is a series of organs located outside the body and around the pelvis region of a male that contribute towards the reproduction process. The primary direct function of the male reproductive system is to provide the male sperm for fertilization of the ovum.
The major reproductive organs of the male can be grouped into three categories. The first category produces and stores sperm (spermatozoa). These are produced in the testicles, which are housed in the temperature-regulating scrotum; immature sperm then travel to the epididymides for development and storage. The second category are the ejaculatory fluid producing glands which include the Cowper's gland (also called bulbourethral gland), seminal vesicles, prostate, and vas deferens. The final category are those used for copulation and deposition of the sperm within the female; these include the penis, urethra, and vas deferens.
Major secondary sexual characteristics include a larger, more muscular stature, deepened voice, facial and body hair, broad shoulders, and the development of an Adam's apple. An important sexual hormone of males is androgen, particularly testosterone.
The testes release a hormone that controls the development of sperm. This hormone is also responsible for the development of physical characteristics in men, such as facial hair and a deep voice.
Development
The development of the reproductive system and the development of the urinary system are closely tied to the development of the human fetus. Despite the differences between them, the adult male and female are determined in early development in the 6th week. The gonads and external genitals are derived from the intermediate mesoderm. The three main fetal precursors of the reproductive organs are the Wolffian duct, the Müllerian ducts, and the gonads. Endocrine hormones are a well-known and critical controlling factor in the normal differentiation of the reproductive system.
The Wolffian duct forms the epididymis, vas deferens, ejaculatory duct, and seminal vesicle in the male reproductive system, but essentially disappears in the female reproductive system. The reverse is true for the Müllerian duct, as it essentially disappears in the male reproductive system and forms the fallopian tubes, uterus, and vagina in the female system. In both sexes, the gonads go on to form the testes and ovaries; because they are derived from the same undeveloped structure, they are considered homologous organs. There are a number of other homologous structures shared between male and female reproductive systems. However, despite the similarity in function of the female fallopian tubes and the male epididymis and vas deferens, they are not homologous but rather analogous structures as they arise from different fetal structures.
Reproduction
Production of gametes
Gametes are produced within the gonads through a process known as gametogenesis. This occurs when certain types of germ cells undergo meiosis to split the normal diploid number of chromosomes (n=46) into haploid cells containing only 23 chromosomes.
In males, this process is known as spermatogenesis and occurs only after puberty in the seminiferous tubules of the testes. The immature spermatozoa or sperm are then sent to the epididymis, where they gain a tail, enabling motility. Each of the original diploid germ cells or primary spermatocytes forms four functional gametes, each forever young. The production and survival of sperms require a temperature below the normal core body temperature. Since the scrotum, where the testes is present, is situated outside the body cavity, it provides a temperature about 3 °C below normal body temperature.
In females, gametogenesis is known as oogenesis; this occurs in the ovarian follicles of the ovaries. This process does not produce mature ovum until puberty. In contrast with males, each of the original diploid germ cells or primary oocytes will form only one mature ovum, and three polar bodies which are not capable of fertilization. It has long been understood that in females, unlike males, all the primary oocytes ever found in a female will be created prior to birth, and that the final stages of ova production will then not resume until puberty. However, recent scientific research has challenged that hypothesis. This new research indicates that in at least some species of mammal, oocytes continue to be replenished in females well after birth.
In male germ cells and spermatozoa, and also in female oocytes, special DNA repair mechanism are present that function to maintain the integrity of the genomes that are to be passed on to progeny. These DNA repair pathways include homologous recombinational repair, non-homologous end joining, base excision repair and DNA mismatch repair.
Disease
Like all complex organ systems, the human reproductive system is affected by many diseases. There are four main categories of reproductive diseases in humans. They are:
Genetic or congenital abnormalities.
Cancers.
Infections, which are often sexually transmitted infections.
Varicocele
Functional problems caused by environmental factors, physical damage, psychological issues, autoimmune disorders, or other causes. The best-known functional problems include sexual dysfunction and infertility, which are both broad terms relating to many disorders with many causes.
Specific reproductive diseases are often symptoms of other diseases and disorders, or have multiple, or unknown causes making them difficult to classify. Examples of unclassifiable disorders are Peyronie's disease in males and endometriosis in females. Many congenital conditions cause reproductive abnormalities, but are better known for their other symptoms. These include: Turner syndrome, Klinefelter's syndrome, cystic fibrosis, and Bloom syndrome.
| Biology and health sciences | Human anatomy | Health |
1906321 | https://en.wikipedia.org/wiki/Ransomware | Ransomware | Ransomware is a type of malware that encrypts the victim's personal data until a ransom is paid. They commonly use difficult-to-trace digital currencies such as paysafecard or Bitcoin and other cryptocurrencies are used for the ransoms, making tracing and prosecuting the perpetrators difficult. Sometimes the original files can be retrieved without paying the ransom due to implementation mistakes, leaked cryptographic keys or a complete lack of encryption in the ransomware.
Ransomware attacks are typically carried out using a Trojan disguised as a legitimate file that the user is tricked into downloading or opening when it arrives as an email attachment. However, one high-profile example, the WannaCry worm, traveled automatically between computers without user interaction.
Starting as early as 1989 with the first documented ransomware known as the AIDS trojan, the use of ransomware scams has grown internationally. There were 181.5 million ransomware attacks in the first six months of 2018. This record marks a 229% increase over this same time frame in 2017. In June 2014, vendor McAfee released data showing that it had collected more than double the number of ransomware samples that quarter than it had in the same quarter the previous year. CryptoLocker was particularly successful, procuring an estimated US$3 million before it was taken down by authorities, and CryptoWall was estimated by the US Federal Bureau of Investigation (FBI) to have accrued over US$18 million by June 2015. In 2020, the IC3 received 2,474 complaints identified as ransomware with adjusted losses of over $29.1 million. The losses could exceed this amount, according to the FBI. Globally, according to Statistica, there were about 623 million ransomware attacks in 2021, and 493 million in 2022.
Operation
The concept of file-encrypting ransomware was invented and implemented by Young and Yung at Columbia University and was presented at the 1996 IEEE Security & Privacy conference. It is called cryptoviral extortion and it was inspired by the fictional facehugger in the movie Alien. Cryptoviral extortion is the following three-round protocol carried out between the attacker and the victim.
[attacker→victim] The attacker generates a key pair and places the corresponding public key in the malware. The malware is released.
[victim→attacker] To carry out the cryptoviral extortion attack, the malware generates a random symmetric key and encrypts the victim's data with it. It uses the public key in the malware to encrypt the symmetric key. This is known as hybrid encryption and it results in a small asymmetric ciphertext as well as the symmetric ciphertext of the victim's data. It zeroizes the symmetric key and the original plaintext data to prevent recovery. It puts up a message to the user that includes the asymmetric ciphertext and how to pay the ransom. The victim sends the asymmetric ciphertext and e-money to the attacker.
[attacker→victim] The attacker receives the payment, deciphers the asymmetric ciphertext with the attacker's private key, and sends the symmetric key to the victim. The victim deciphers the encrypted data with the needed symmetric key thereby completing the cryptovirology attack.
The symmetric key is randomly generated and will not assist other victims. At no point is the attacker's private key exposed to victims and the victim need only send a very small ciphertext (the encrypted symmetric-cipher key) to the attacker.
Ransomware attacks are typically carried out using a Trojan, entering a system through, for example, a malicious attachment, an embedded link in a phishing email, or a vulnerability in a network service. The program then runs a payload, which locks the system in some fashion, or claims to lock the system but does not (e.g., a scareware program). Payloads may display a fake warning purportedly by an entity such as a law enforcement agency, falsely claiming that the system has been used for illegal activities, contains content such as pornography and "pirated" media.
Some payloads consist simply of an application designed to lock or restrict the system until payment is made, typically by setting the Windows Shell to itself, or even modifying the master boot record and/or partition table to prevent the operating system from booting until it is repaired. The most sophisticated payloads encrypt files, with many using strong encryption to encrypt the victim's files in such a way that only the malware author has the needed decryption key.
Payment is virtually always the goal, and the victim is coerced into paying for the ransomware to be removed either by supplying a program that can decrypt the files, or by sending an unlock code that undoes the payload's changes. While the attacker may simply take the money without returning the victim's files, it is in the attacker's best interest to perform the decryption as agreed, since victims will stop sending payments if it becomes known that they serve no purpose. A key element in making ransomware work for the attacker is a convenient payment system that is hard to trace. A range of such payment methods have been used, including wire transfers, premium-rate text messages, pre-paid voucher services such as paysafecard, and the Bitcoin cryptocurrency.
In May 2020, vendor Sophos reported that the global average cost to remediate a ransomware attack (considering downtime, people time, device cost, network cost, lost opportunity and ransom paid) was $761,106. Ninety-five percent of organizations that paid the ransom had their data restored.
History
Encrypting ransomware
The first known malware extortion attack, the "AIDS Trojan" written by Joseph Popp in 1989, had a design failure so severe it was not necessary to pay the extortionist at all. Its payload hid the files on the hard drive and encrypted only their names, and displayed a message claiming that the user's license to use a certain piece of software had expired. The user was asked to pay US$189 to "PC Cyborg Corporation" in order to obtain a repair tool even though the decryption key could be extracted from the code of the Trojan. The Trojan was also known as "PC Cyborg". Popp was declared mentally unfit to stand trial for his actions, but he promised to donate the profits from the malware to fund AIDS research.
The idea of abusing anonymous cash systems to safely collect ransom from human kidnapping was introduced in 1992 by Sebastiaan von Solms and David Naccache. This electronic money collection method was also proposed for cryptoviral extortion attacks. In the von Solms-Naccache scenario a newspaper publication was used (since bitcoin ledgers did not exist at the time the paper was written).
The notion of using public key cryptography for data kidnapping attacks was introduced in 1996 by Adam L. Young and Moti Yung. Young and Yung critiqued the failed AIDS Information Trojan that relied on symmetric cryptography alone, the fatal flaw being that the decryption key could be extracted from the Trojan, and implemented an experimental proof-of-concept cryptovirus on a Macintosh SE/30 that used RSA and the Tiny Encryption Algorithm (TEA) to hybrid encrypt the victim's data. Since public key cryptography is used, the virus only contains the encryption key. The attacker keeps the corresponding private decryption key private. Young and Yung's original experimental cryptovirus had the victim send the asymmetric ciphertext to the attacker who deciphers it and returns the symmetric decryption key it contains to the victim for a fee. Long before electronic money existed Young and Yung proposed that electronic money could be extorted through encryption as well, stating that "the virus writer can effectively hold all of the money ransom until half of it is given to him. Even if the e-money was previously encrypted by the user, it is of no use to the user if it gets encrypted by a cryptovirus". They referred to these attacks as being "cryptoviral extortion", an overt attack that is part of a larger class of attacks in a field called cryptovirology, which encompasses both overt and covert attacks. The cryptoviral extortion protocol was inspired by the parasitic relationship between H. R. Giger's facehugger and its host in the movie Alien.
Examples of extortionate ransomware became prominent in May 2005. By mid-2006, Trojans such as Gpcode, TROJ.RANSOM.A, Archiveus, Krotten, Cryzip, and MayArchive began utilizing more sophisticated RSA encryption schemes, with ever-increasing key-sizes. Gpcode.AG, which was detected in June 2006, was encrypted with a 660-bit RSA public key. In June 2008, a variant known as Gpcode.AK was detected. Using a 1024-bit RSA key, it was believed large enough to be computationally infeasible to break without a concerted distributed effort.
Encrypting ransomware returned to prominence in late 2013 with the propagation of CryptoLocker—using the Bitcoin digital currency platform to collect ransom money. In December 2013, ZDNet estimated based on Bitcoin transaction information that between 15 October and 18 December, the operators of CryptoLocker had procured about US$27 million from infected users. The CryptoLocker technique was widely copied in the months following, including CryptoLocker 2.0 (thought not to be related to CryptoLocker), CryptoDefense (which initially contained a major design flaw that stored the private key on the infected system in a user-retrievable location, due to its use of Windows' built-in encryption APIs), and the August 2014 discovery of a Trojan specifically targeting network-attached storage devices produced by Synology. In January 2015, it was reported that ransomware-styled attacks have occurred against individual websites via hacking, and through ransomware designed to target Linux-based web servers.
In 2022, Costa Rica received widespread Conti ransomware attacks affecting government, healthcare and industry. This led President Rodrigo Chaves to declare a state of emergency and announce that Costa Rica is "at war" with its ransomware hackers.
In some infections, there is a two-stage payload, common in many malware systems. The user is tricked into running a script, which downloads the main virus and executes it. In early versions of the dual-payload system, the script was contained in a Microsoft Office document with an attached VBScript macro, or in a windows scripting facility (WSF) file. As detection systems started blocking these first stage payloads, the Microsoft Malware Protection Center identified a trend away toward LNK files with self-contained Microsoft Windows PowerShell scripts. In 2016, PowerShell was found to be involved in nearly 40% of endpoint security incidents.
Some ransomware strains have used proxies tied to Tor hidden services to connect to their command and control servers, increasing the difficulty of tracing the exact location of the criminals. Furthermore, dark web vendors have increasingly started to offer the technology as a service, wherein ransomware is sold, ready for deployment on victims' machines, on a subscription basis, similarly to Adobe Creative Cloud or Office 365.
Symantec has classified ransomware to be the most dangerous cyber threat.
Non-encrypting ransomware
In August 2010, Russian authorities arrested nine individuals connected to a ransomware Trojan known as WinLock. Unlike the previous Gpcode Trojan, WinLock did not use encryption. Instead, WinLock trivially restricted access to the system by displaying pornographic images and asked users to send a premium-rate SMS (costing around US$10) to receive a code that could be used to unlock their machines. The scam hit numerous users across Russia and neighbouring countries—reportedly earning the group over US$16 million.
In 2011, a ransomware Trojan surfaced that imitated the Windows Product Activation notice, and informed users that a system's Windows installation had to be re-activated due to "[being a] victim of fraud". An online activation option was offered (like the actual Windows activation process), but was unavailable, requiring the user to call one of six international numbers to input a 6-digit code. While the malware claimed that this call would be free, it was routed through a rogue operator in a country with high international phone rates, who placed the call on hold, causing the user to incur large international long-distance charges.
In 2012, Symantec reported spread out of Eastern Europe of ransomware with a lock screen purporting to be law enforcement demanding payment for illegal activity.
In February 2013, a ransomware Trojan based on the Stamp.EK exploit kit surfaced; the malware was distributed via sites hosted on the project hosting services SourceForge and GitHub that claimed to offer "fake nude pics" of celebrities. In July 2013, an OS X-specific ransomware Trojan surfaced, which displays a web page that accuses the user of downloading pornography. Unlike its Windows-based counterparts, it does not block the entire computer, but simply exploits the behaviour of the web browser itself to frustrate attempts to close the page through normal means.
In July 2013, a 21-year-old man from Virginia, whose computer coincidentally did contain pornographic photographs of underage girls with whom he had conducted sexualized communications, turned himself in to police after receiving and being deceived by FBI MoneyPak Ransomware accusing him of possessing child pornography. An investigation discovered the incriminating files, and the man was charged with child sexual abuse and possession of child pornography.
Exfiltration (Leakware / Doxware)
The converse of ransomware is a cryptovirology attack invented by Adam L. Young that threatens to publish stolen information from the victim's computer system rather than deny the victim access to it. In a leakware attack, malware exfiltrates sensitive host data either to the attacker or alternatively, to remote instances of the malware, and the attacker threatens to publish the victim's data unless a ransom is paid. The attack was presented at West Point in 2003 and was summarized in the book Malicious Cryptography as follows, "The attack differs from the extortion attack in the following way. In the extortion attack, the victim is denied access to its own valuable information and has to pay to get it back, where in the attack that is presented here the victim retains access to the information but its disclosure is at the discretion of the computer virus". The attack is rooted in game theory and was originally dubbed "non-zero sum games and survivable malware". The attack can yield monetary gain in cases where the malware acquires access to information that may damage the victim user or organization, e.g., the reputational damage that could result from publishing proof that the attack itself was a success.
Common targets for exfiltration include:
third party information stored by the primary victim (such as customer account information or health records);
information proprietary to the victim (such as trade secrets and product information)
embarrassing information (such as the victim's health information or information about the victim's personal past)
Exfiltration attacks are usually targeted, with a curated victim list, and often preliminary surveillance of the victim's systems to find potential data targets and weaknesses.
Mobile ransomware
With the increased popularity of ransomware on PC platforms, ransomware targeting mobile operating systems has also proliferated. Typically, mobile ransomware payloads are blockers, as there is little incentive to encrypt data since it can be easily restored via online synchronization. Mobile ransomware typically targets the Android platform, as it allows applications to be installed from third-party sources. The payload is typically distributed as an APK file installed by an unsuspecting user; it may attempt to display a blocking message over top of all other applications, while another used a form of clickjacking to cause the user to give it "device administrator" privileges to achieve deeper access to the system.
Different tactics have been used on iOS devices, such as exploiting iCloud accounts and using the Find My iPhone system to lock access to the device. On iOS 10.3, Apple patched a bug in the handling of JavaScript pop-up windows in Safari that had been exploited by ransomware websites.
It recently has been shown that ransomware may also target ARM architectures like those that can be found in various Internet-of-Things (IoT) devices, such as Industrial IoT edge devices.
In August 2019 researchers demonstrated it is possible to infect DSLR cameras with ransomware. Digital cameras often use Picture Transfer Protocol (PTP - standard protocol used to transfer files.) Researchers found that it was possible to exploit vulnerabilities in the protocol to infect target camera(s) with ransomware (or execute any arbitrary code). This attack was presented at the Defcon security conference in Las Vegas as a proof of concept attack (not as actual armed malware).
Progression of attacks
The first attacks were on random users, typically infected through email attachments sent by small groups of criminals, demanding a few hundred dollars in cryptocurrency to unlock files (typically a private individual's photographs and documents) that the ransomware had encrypted. As ransomware matured as a business, organised gangs entered the field, advertising on the dark Web for experts, and outsourcing functions. This led to improvement in the quality of ransomware and its success. Rather than random emails, the gangs stole credentials, found vulnerabilities in target networks, and improved the malware to avoid detection by anti-malware scanners. Ransoms demanded escalated into the much larger sums (millions) that an enterprise would pay to recover its data, rather than what an individual would pay for their documents (hundreds).
In 2016, a significant uptick in ransomware attacks on hospitals was noted. According to the 2017 Internet Security Threat Report from Symantec Corp, ransomware affected not only IT systems but also patient care, clinical operations, and billing. Online criminals may be motivated by the money available and sense of urgency within the healthcare system.
Ransomware is growing rapidly across the internet users but also for the IoT environment. The big problem is that millions of dollars are lost by some organizations and industries that have decided to pay, such as the Hollywood Presbyterian Medical Center and the MedStar Health.
According to Symantec 2019 ISTR report, for the first time since 2013, in 2018 there was an observed decrease in ransomware activity with a drop of 20 percent. Before 2017, consumers were the preferred victims, but in 2017 this changed dramatically, it moved to the enterprises. In 2018 this path accelerated with 81 percent infections which represented a 12 percent increase. The common distribution method today is based on email campaigns.
In late 2019 ransomware group Maze downloaded companies' sensitive files before locking them, and threatened to leak the data publicly if the ransom was not paid; in at least one case they did this. Many other gangs followed; "leak sites" were created on the dark web where stolen data could be accessed. Later attacks focussed on the threat to leak data, without necessarily locking it—this negated the protection afforded victims by robust backup procedures. there is a risk of hostile governments using ransomware to conceal what is actually intelligence gathering.
The first reported death following a ransomware attack was at a German hospital in October 2020.
A significant increase in ransomware attacks occurred during the 2020 COVID-19 pandemic. Evidence has demonstrated that the targeted institutions of these attacks included government, finance, and healthcare. Researchers have contended that several different factors can explain the increase in attacks during this time. However, a major factor is that remote work, which became the norm for many industries in 2020, led to the surge in attacks because of the lack of security in comparison to traditional work environments.
Notable attack targets
Notable software packages
Reveton
In 2012, a major ransomware Trojan known as Reveton began to spread. Based on the Citadel Trojan (which, itself, is based on the Zeus Trojan), its payload displays a warning purportedly from a law enforcement agency claiming that the computer has been used for illegal activities, such as downloading unlicensed software or child pornography. Due to this behaviour, it is commonly referred to as the "Police Trojan". The warning informs the user that to unlock their system, they would have to pay a fine using a voucher from an anonymous prepaid cash service such as Ukash or paysafecard. To increase the illusion that the computer is being tracked by law enforcement, the screen also displays the computer's IP address, while some versions display footage from a victim's webcam to give the illusion that the user is being recorded.
Reveton initially began spreading in various European countries in early 2012. Variants were localized with templates branded with the logos of different law enforcement organizations based on the user's country; for example, variants used in the United Kingdom contained the branding of organizations such as the Metropolitan Police Service and the Police National E-Crime Unit. Another version contained the logo of the royalty collection society PRS for Music, which specifically accused the user of illegally downloading music. In a statement warning the public about the malware, the Metropolitan Police clarified that they would never lock a computer in such a way as part of an investigation.
In May 2012, Trend Micro threat researchers discovered templates for variations for the United States and Canada, suggesting that its authors may have been planning to target users in North America. By August 2012, a new variant of Reveton began to spread in the United States, claiming to require the payment of a $200 fine to the FBI using a MoneyPak card. In February 2013, a Russian citizen was arrested in Dubai by Spanish authorities for his connection to a crime ring that had been using Reveton; ten other individuals were arrested on money laundering charges. In August 2014, Avast Software reported that it had found new variants of Reveton that also distribute password-stealing malware as part of its payload.
CryptoLocker
Encrypting ransomware reappeared in September 2013 with a Trojan known as CryptoLocker, which generated a 2048-bit RSA key pair and uploaded in turn to a command-and-control server, and used to encrypt files using a whitelist of specific file extensions. The malware threatened to delete the private key if a payment of Bitcoin or a pre-paid cash voucher was not made within 3 days of the infection. Due to the extremely large key size it uses, analysts and those affected by the Trojan considered CryptoLocker extremely difficult to repair. Even after the deadline passed, the private key could still be obtained using an online tool, but the price would increase to 10 BTC—which cost approximately US$2300 as of November 2013.
CryptoLocker was isolated by the seizure of the Gameover ZeuS botnet as part of Operation Tovar, as officially announced by the U.S. Department of Justice on 2 June 2014. The Department of Justice also publicly issued an indictment against the Russian hacker Evgeniy Bogachev for his alleged involvement in the botnet. It was estimated that at least US$3 million was extorted with the malware before the shutdown.
CryptoLocker.F and TorrentLocker
In September 2014, a wave of ransomware Trojans surfaced that first targeted users in Australia, under the names CryptoWall and CryptoLocker (which is, as with CryptoLocker 2.0, unrelated to the original CryptoLocker). The Trojans spread via fraudulent e-mails claiming to be failed parcel delivery notices from Australia Post; to evade detection by automatic e-mail scanners that follow all links on a page to scan for malware, this variant was designed to require users to visit a web page and enter a CAPTCHA code before the payload is actually downloaded, preventing such automated processes from being able to scan the payload. Symantec determined that these new variants, which it identified as CryptoLocker.F, were again, unrelated to the original CryptoLocker due to differences in their operation. A notable victim of the Trojans was the Australian Broadcasting Corporation; live programming on its television news channel ABC News 24 was disrupted for half an hour and shifted to Melbourne studios due to a CryptoWall infection on computers at its Sydney studio.
Another Trojan in this wave, TorrentLocker, initially contained a design flaw comparable to CryptoDefense; it used the same keystream for every infected computer, making the encryption trivial to overcome. However, this flaw was later fixed. By late-November 2014, it was estimated that over 9,000 users had been infected by TorrentLocker in Australia alone, trailing only Turkey with 11,700 infections.
CryptoWall
Another major ransomware Trojan targeting Windows, CryptoWall, first appeared in 2014. One strain of CryptoWall was distributed as part of a malvertising campaign on the Zedo ad network in late-September 2014 that targeted several major websites; the ads redirected to rogue websites that used browser plugin exploits to download the payload. A Barracuda Networks researcher also noted that the payload was signed with a digital signature in an effort to appear trustworthy to security software.
CryptoWall 3.0 used a payload written in JavaScript as part of an email attachment, which downloads executables disguised as JPG images. To further evade detection, the malware creates new instances of explorer.exe and svchost.exe to communicate with its servers. When encrypting files, the malware also deletes volume shadow copies and installs spyware that steals passwords and Bitcoin wallets.
The FBI reported in June 2015 that nearly 1,000 victims had contacted the bureau's Internet Crime Complaint Center to report CryptoWall infections, and estimated losses of at least $18 million.
The most recent version, CryptoWall 4.0, enhanced its code to avoid antivirus detection, and encrypts not only the data in files but also the file names.
Fusob
Fusob is a major family of mobile ransomware. Between April 2015 and March 2016, about 56 percent of accounted mobile ransomware was Fusob.
Like most other pieces of ransomware, it employs scare tactics to extort a hefty sum from the user. The app acts as if it were a notice from the authorities, demanding the victim to pay a fine from $100 to $200 USD or otherwise face a fictitious criminal charge. Fusob requests iTunes gift cards for payment, unlike most cryptocurrency-centric ransomware.
In order to infect devices, Fusob masquerades as a pornographic video player. When it is installed, it first checks the device's system language. If the language is Russian or Eastern-European, Fusob remains dormant. Otherwise, it locks the device and demands ransom. About 40% of victims are in Germany, while the United Kingdom encompasses 14.5% of victims and the US encompasses 11.4%. Fusob and Small (another family of ransomware) represented over 93% of mobile ransomware between 2015 and 2016.
WannaCry
In May 2017, the WannaCry ransomware attack spread through the Internet, using an exploit vector named EternalBlue, which was allegedly leaked from the U.S. National Security Agency. The ransomware attack, unprecedented in scale, infected more than 230,000 computers in over 150 countries, using 20 different languages to demand money from users using Bitcoin cryptocurrency. WannaCry demanded US$300 per computer. The attack affected Telefónica and several other large companies in Spain, as well as parts of the British National Health Service (NHS), where at least 16 hospitals had to turn away patients or cancel scheduled operations, FedEx, Deutsche Bahn, Honda, Renault, as well as the Russian Interior Ministry and Russian telecom MegaFon. The attackers gave their victims a 7-day deadline from the day their computers got infected, after which the encrypted files would be deleted.
Petya
Petya was first discovered in March 2016; unlike other forms of encrypting ransomware, the malware aimed to infect the master boot record, installing a payload which encrypts the file tables of the NTFS file system the next time that the infected system boots, blocking the system from booting into Windows at all until the ransom is paid. Check Point reported that despite what it believed to be an innovative evolution in ransomware design, it had resulted in relatively-fewer infections than other ransomware active around the same time frame.
On 27 June 2017, a heavily modified version of Petya was used for a global cyberattack primarily targeting Ukraine (but affecting many countries). This version had been modified to propagate using the same EternalBlue exploit that was used by WannaCry. Due to another design change, it is also unable to actually unlock a system after the ransom is paid; this led to security analysts speculating that the attack was not meant to generate illicit profit, but to simply cause disruption.
Bad Rabbit
On 24 October 2017, some users in Russia and Ukraine reported a new ransomware attack, named "Bad Rabbit", which follows a similar pattern to WannaCry and Petya by encrypting the user's file tables and then demands a Bitcoin payment to decrypt them. ESET believed the ransomware to have been distributed by a bogus update to Adobe Flash software. Among agencies that were affected by the ransomware were: Interfax, Odesa International Airport, Kyiv Metro, and the Ministry of Infrastructure of Ukraine. As it used corporate network structures to spread, the ransomware was also discovered in other countries, including Turkey, Germany, Poland, Japan, South Korea, and the United States. Experts believed the ransomware attack was tied to the Petya attack in Ukraine (especially because Bad Rabbit's code has many overlapping and analogical elements to the code of Petya/NotPetya, appending to CrowdStrike Bad Rabbit and NotPetya's dynamic link library (DLL) share 67 percent of the same code) though the only identity to the culprits are the names of characters from the Game of Thrones series embedded within the code.
Security experts found that the ransomware did not use the EternalBlue exploit to spread, and a simple method to inoculate an unaffected machine running older Windows versions was found by 24 October 2017. Further, the sites that had been used to spread the bogus Flash updating have gone offline or removed the problematic files within a few days of its discovery, effectively killing off the spread of Bad Rabbit.
SamSam
In 2016, a new strain of ransomware emerged that was targeting JBoss servers. This strain, named "SamSam", was found to bypass the process of phishing or illicit downloads in favor of exploiting vulnerabilities on weak servers. The malware uses a Remote Desktop Protocol brute-force attack to guess weak passwords until one is broken. The virus has been behind attacks on government and healthcare targets, with notable hacks occurring against the town of Farmington, New Mexico, the Colorado Department of Transportation, Davidson County, North Carolina, and most recently, a ransomware attack on the infrastructure of Atlanta.
Mohammad Mehdi Shah Mansouri (born in Qom, Iran in 1991) and Faramarz Shahi Savandi (born in Shiraz, Iran, in 1984) are wanted by the FBI for allegedly launching SamSam ransomware. The two have allegedly made $6 million from extortion and caused over $30 million in damages using the malware.
DarkSide
On May 7, 2021, a cyberattack was executed on the US Colonial Pipeline. The Federal Bureau of Investigation identified DarkSide as the perpetrator of the Colonial Pipeline ransomware attack, perpetrated by malicious code, that led to a voluntary shutdown of the main pipeline supplying 45% of fuel to the East Coast of the United States. The attack was described as the worst cyberattack to date on the U.S. critical infrastructure. DarkSide successfully extorted about 75 Bitcoin (almost US$5 million) from Colonial Pipeline. U.S. officials are investigating whether the attack was purely criminal or took place with the involvement of the Russian government or another state sponsor. Following the attack, DarkSide posted a statement claiming that "We are apolitical, we do not participate in geopolitics...Our goal is to make money and not creating problems for society."
In May 2021, the FBI and Cybersecurity and Infrastructure Security Agency (CISA) issued a joint alert urging the owners and operators of critical infrastructure to take certain steps to reduce their vulnerability to DarkSide ransomware and ransomware in general.
Syskey
Syskey is a utility that was included with Windows NT-based operating systems to encrypt the user account database, optionally with a password. The tool has sometimes been effectively used as ransomware during technical support scams—where a caller with remote access to the computer may use the tool to lock the user out of their computer with a password known only to them. Syskey was removed from later versions of Windows 10 and Windows Server in 2017, due to being obsolete and "known to be used by hackers as part of ransomware scams".
Ransomware-as-a-service
Ransomware-as-a-service (RaaS) became a notable method after the Russia-based or Russian-speaking group REvil staged operations against several targets, including the Brazil-based JBS S.A. in May 2021, and the US-based Kaseya Limited in July 2021. After a July 9, 2021 phone call between United States president Joe Biden and Russian president Vladimir Putin, Biden told the press, "I made it very clear to him that the United States expects when a ransomware operation is coming from his soil even though it’s not sponsored by the state, we expect them to act if we give them enough information to act on who that is." Biden later added that the United States would take the group's servers down if Putin did not. Four days later, REvil websites and other infrastructure vanished from the internet.
Mitigation
If an attack is suspected or detected in its early stages, it takes some time for encryption to take place; immediate removal of the malware (a relatively simple process) before it has completed would stop further damage to data, without salvaging any already lost.
Security experts have suggested precautionary measures for dealing with ransomware. Using software or other security policies to block known payloads from launching will help to prevent infection, but will not protect against all attacks As such, having a proper backup solution is a critical component to defending against ransomware. Note that, because many ransomware attackers will not only encrypt the victim's live machine but it will also attempt to delete any hot backups stored locally or on accessible over the network on a NAS, it is also critical to maintain "offline" backups of data stored in locations inaccessible from any potentially infected computer, such as external storage drives or devices that do not have any access to any network (including the Internet), prevents them from being accessed by the ransomware. Moreover, if using a NAS or Cloud storage, then the computer should have append-only permission to the destination storage, such that it cannot delete or overwrite previous backups. According to comodo, applying two Attack Surface Reduction on OS/Kernel provides a materially-reduced attack surface which results in a heightened security posture.
Installing security updates issued by software vendors can mitigate the vulnerabilities leveraged by certain strains to propagate. Other measures include cyber hygiene − exercising caution when opening e-mail attachments and links, network segmentation, and keeping critical computers isolated from networks. Furthermore, to mitigate the spread of ransomware measures of infection control can be applied. Such may include disconnecting infected machines from all networks, educational programs, effective communication channels, malware surveillance and ways of collective participation
In August 2021, the Cybersecurity and Infrastructure Security Agency (CISA) released a report that provided guidance for how to mitigate ransomware attacks. This was due to a significant jump in recent attacks related to ransomware. These attacks included aggression against a US pipeline company and a software company, which impacted the downstream customers of MSPs.
Write Once Read Many (WORM) storage, such as many optical disc formats are virtually immune to ransomware since its contents cannot be changed or deleted. However data cannot be deleted in any way making it impractical for many storage due to privacy laws and other content laws, the only way is to copy it to a new WORM disk minus the unwanted files, then destroy the original copy.
File system defenses against ransomware
A number of file systems keep snapshots of the data they hold, which can be used to recover the contents of files from a time prior to the ransomware attack in the event the ransomware does not disable it.
On Windows, the Volume shadow copy (VSS) is often used to store backups of data; ransomware often targets these snapshots to prevent recovery and therefore it is often advisable to disable user access to the user tool VSSadmin.exe to reduce the risk that ransomware can disable or delete past copies.
On Windows 10, users can add specific directories or files to Controlled Folder Access in Windows Defender to protect them from ransomware. It is advised to add backup and other important directories to Controlled Folder Access.
Unless malware gains root on the ZFS host system in deploying an attack coded to issue ZFS administrative commands, file servers running ZFS are broadly immune to ransomware, because ZFS is capable of snapshotting even a large file system many times an hour, and these snapshots are immutable (read only) and easily rolled back or files recovered in the event of data corruption. In general, only an administrator can delete (but cannot modify) snapshots.
File decryption and recovery
There are a number of tools intended specifically to decrypt files locked by ransomware, although successful recovery may not be possible. If the same encryption key is used for all files, decryption tools use files for which there are both uncorrupted backups and encrypted copies (a known-plaintext attack in the jargon of cryptanalysis. But it only works when the cipher the attacker used was weak to begin with, being vulnerable to known-plaintext attack); recovery of the key, if it is possible, may take several days. Free ransomware decryption tools can help decrypt files encrypted by the following forms of ransomware: AES_NI, Alcatraz Locker, Apocalypse, BadBlock, Bart, BTCWare, Crypt888, CryptoMix, CrySiS, EncrypTile, FindZip, Globe, Hidden Tear, Jigsaw, LambdaLocker, Legion, NoobCrypt, Stampado, SZFLocker, TeslaCrypt, XData. Ransomware encryption that has been cracked by security researchers is typically abandoned for criminal purposes; thus in practice most attacks cannot be reverted by breaking encryption.
The No More Ransom Project is an initiative by the Netherlands' police's National High Tech Crime Unit, Europol’s European Cybercrime Centre, Kaspersky Lab and McAfee to help ransomware victims recover their data without paying a ransom. They offer a free CryptoSheriff tool to analyze encrypted files and search for decryption tools.
In addition, old copies of files may exist on the disk, which has been previously deleted. In some cases, these deleted versions may still be recoverable using software designed for that purpose.
A 2019 ProPublica investigation found the cybersecurity firms Proven Data Recovery and Monstercloud, which advertised ransom-free decryption services, would typically simply pay the ransom and charge the victim a higher price. SamSam hackers dealt with Proven Data so frequently that they would recommend the company to victims having technical difficulties making payment. Other companies like Coveware were more transparent in offering the service of paying the hackers and patching insecure systems. Many American victims found the ransom amount was too low to meet the United States Department of Justice threshold for federal involvement, but that local police lacked the technical capabilities to help and were often victims themselves.
Criminal arrests and convictions
Zain Qaiser
A British student, Zain Qaiser, from Barking, London was jailed for more than six years at Kingston upon Thames Crown Court for his ransomware attacks in 2019. He is said to have been "the most prolific cyber criminal to be sentenced in the UK". He became active when he was only 17. He contacted the Russian controller of one of the most powerful attacks, believed to be the Lurk malware gang, and arranged for a split of his profits. He also contacted online criminals from China and the US to move the money. For about one and a half years, he posed as a legitimate supplier of online promotions of book advertising on some of the world's most visited legal pornography websites. Each of the adverts that were promoted on the websites contained the Reveton Ransomware strain of the malicious Angler Exploit Kit (AEK) that seized control of the machine. Investigators discovered about £700,000 of earnings, although his network may have earned more than £4m. He may have hidden some money using cryptocurrencies. The ransomware would instruct victims to buy GreenDot MoneyPak vouchers and enter the code in the Reveton panel displayed on the screen. This money entered a MoneyPak account managed by Qaiser, who would then deposit the voucher payments into the debit card account of his American co-conspirator, Raymond Odigie Uadiale. Uadiale was a student at Florida International University during 2012 and 2013 and later worked for Microsoft. Uadiale would convert the money into Liberty Reserve digital currency and deposit it into Qaiser's Liberty Reserve account.
A breakthrough, in this case, occurred in May 2013 when authorities from several countries seized the Liberty Reserve servers, obtaining access to all its transactions and account history. Qaiser was running encrypted virtual machines on his Macbook Pro with both Mac and Windows operating systems. He could not be tried earlier because he was sectioned (involuntarily committed) under the UK Mental Health Act of 1983 at Goodmayes Hospital where he was found to be using the hospital Wi-Fi to access his advertising sites. His lawyer claimed that Qaiser had suffered from mental illness. Russian police arrested 50 members of the Lurk malware gang in June 2016. Uadiale, a naturalized US citizen of Nigerian descent, was jailed for 18 months.
Legal aspects
The publication of proof-of-concept attack code is common among academic researchers and vulnerability researchers.
It teaches the nature of the threat, conveys the gravity of the issues, and enables countermeasures to be devised and put into place. However, lawmakers with the support of law-enforcement bodies are contemplating making the creation of ransomware illegal. In the state of Maryland, the original draft of HB 340 made it a felony to create ransomware, punishable by up to 10 years in prison. A minor in Japan was arrested for creating and distributing ransomware code.
| Technology | Computer security | null |
1906854 | https://en.wikipedia.org/wiki/R/K%20selection%20theory | R/K selection theory | In ecology, selection theory relates to the selection of combinations of traits in an organism that trade off between quantity and quality of offspring. The focus on either an increased quantity of offspring at the expense of reduced individual parental investment of -strategists, or on a reduced quantity of offspring with a corresponding increased parental investment of -strategists, varies widely, seemingly to promote success in particular environments. The concepts of quantity or quality offspring are sometimes referred to as "cheap" or "expensive", a comment on the expendable nature of the offspring and parental commitment made. The stability of the environment can predict if many expendable offspring are made or if fewer offspring of higher quality would lead to higher reproductive success. An unstable environment would encourage the parent to make many offspring, because the likelihood of all (or the majority) of them surviving to adulthood is slim. In contrast, more stable environments allow parents to confidently invest in one offspring because they are more likely to survive to adulthood.
The terminology of -selection was coined by the ecologists Robert MacArthur and E. O. Wilson in 1967 based on their work on island biogeography; although the concept of the evolution of life history strategies has a longer history (see e.g. plant strategies).
The theory was popular in the 1970s and 1980s, when it was used as a heuristic device, but lost importance in the early 1990s, when it was criticized by several empirical studies. A life-history paradigm has replaced the selection paradigm, but continues to incorporate its important themes as a subset of life history theory. Some scientists now prefer to use the terms fast versus slow life history as a replacement for, respectively, versus reproductive strategy.
Overview
In selection theory, selective pressures are hypothesised to drive evolution in one of two generalized directions: - or -selection. These terms, and , are drawn from standard ecological formula as illustrated in the simplified Verhulst model of population dynamics:
where is the population, is the maximum growth rate, is the carrying capacity of the local environment, and (the derivative of population size with respect to time ) is the rate of change in population with time. Thus, the equation relates the growth rate of the population to the current population size, incorporating the effect of the two constant parameters and .
(Note that when the population size is greater than the carrying capacity then 1 - N/K is negative, which indicates a population decline or negative growth.) The choice of the letter came from the German Kapazitätsgrenze (capacity limit), while came from rate.
r-selection
-selected species are those that emphasize high growth rates, typically exploit less-crowded ecological niches, and produce many offspring, each of which has a relatively low probability of surviving to adulthood (i.e., high , low ). A typical species is the dandelion (genus Taraxacum).
In unstable or unpredictable environments, -selection predominates due to the ability to reproduce rapidly. There is little advantage in adaptations that permit successful competition with other organisms, because the environment is likely to change again. Among the traits that are thought to characterize -selection are high fecundity, small body size, early maturity onset, short generation time, and the ability to disperse offspring widely.
Organisms whose life history is subject to -selection are often referred to as -strategists or -selected. Organisms that exhibit -selected traits can range from bacteria and diatoms, to insects and grasses, to various semelparous cephalopods, certain families of birds, such as dabbling ducks, and small mammals, particularly rodents.
K-selection
By contrast, -selected species display traits associated with living at densities close to carrying capacity and typically are strong competitors in such crowded niches, that invest more heavily in fewer offspring, each of which has a relatively high probability of surviving to adulthood (i.e., low , high ). In scientific literature, -selected species are occasionally referred to as "opportunistic" whereas -selected species are described as "equilibrium".
In stable or predictable environments, -selection predominates as the ability to compete successfully for limited resources is crucial and populations of -selected organisms typically are very constant in number and close to the maximum that the environment can bear (unlike -selected populations, where population sizes can change much more rapidly).
Traits that are thought to be characteristic of -selection include large body size, long life expectancy, and the production of fewer offspring, which often require extensive parental care until they mature. Organisms whose life history is subject to -selection are often referred to as -strategists or -selected. Organisms with -selected traits include large organisms such as elephants, sharks, humans, and whales, but also smaller long-lived organisms such as Arctic terns, parrots, and eagles.
Continuous spectrum
Although some organisms are identified as primarily - or -strategists, the majority of organisms do not follow this pattern. For instance, trees have traits such as longevity and strong competitiveness that characterise them as -strategists. In reproduction, however, trees typically produce thousands of offspring and disperse them widely, traits characteristic of -strategists.
Similarly, reptiles such as sea turtles display both - and -traits: Although sea turtles are large organisms with long lifespans (provided they reach adulthood), they produce large numbers of unnurtured offspring.
The dichotomy can be re-expressed as a continuous spectrum using the economic concept of discounted future returns, with -selection corresponding to large discount rates and -selection corresponding to small discount rates.
Ecological succession
In areas of major ecological disruption or sterilisation (such as after a major volcanic eruption, as at Krakatoa or Mount St. Helens), - and -strategists play distinct roles in the ecological succession that regenerates the ecosystem. Because of their higher reproductive rates and ecological opportunism, primary colonisers typically are -strategists and they are followed by a succession of increasingly competitive flora and fauna. The ability of an environment to increase energetic content, through photosynthetic capture of solar energy, increases with the increase in complex biodiversity as species proliferate to reach a peak possible with strategies.
Eventually a new equilibrium is approached (sometimes referred to as a climax community), with -strategists gradually being replaced by -strategists which are more competitive and better adapted to the emerging micro-environmental characteristics of the landscape. Traditionally, biodiversity was considered maximized at this stage, with introductions of new species resulting in the replacement and local extinction of endemic species. However, the intermediate disturbance hypothesis posits that intermediate levels of disturbance in a landscape create patches at different levels of succession, promoting coexistence of colonizers and competitors at the regional scale.
Application
While usually applied at the level of species, selection theory is also useful in studying the evolution of ecological and life history differences between subspecies, for instance the African honey bee, A. m. scutellata, and the Italian bee, A. m. ligustica. At the other end of the scale, it has also been used to study the evolutionary ecology of whole groups of organisms, such as bacteriophages. Other researchers have proposed that the evolution of human inflammatory responses is related to selection.
Some researchers, such as Lee Ellis, J. Philippe Rushton, and Aurelio José Figueredo, have attempted to apply selection theory to various human behaviors, including crime, sexual promiscuity, fertility, IQ, and other traits related to life history theory. Rushton developed "differential theory" to attempt to explain variations in behavior across human races. Differential theory has been debunked as being devoid of empirical basis, and has also been described as a key example of scientific racism.
Status
Although selection theory became widely used during the 1970s, it also began to attract more critical attention. In particular, a review in 1977 by the ecologist Stephen C. Stearns drew attention to gaps in the theory, and to ambiguities in the interpretation of empirical data for testing it.
In 1981, a review of the selection literature by Parry demonstrated that there was no agreement among researchers using the theory about the definition of - and -selection, which led him to question whether the assumption of a relation between reproductive expenditure and packaging of offspring was justified. A 1982 study by Templeton and Johnson showed that in a population of Drosophila mercatorum under -selection the population actually produced a higher frequency of traits typically associated with -selection. Several other studies contradicting the predictions of selection theory were also published between 1977 and 1994.
When Stearns reviewed the status of the theory again in 1992, he noted that from 1977 to 1982 there was an average of 42 references to the theory per year in the BIOSIS literature search service, but from 1984 to 1989 the average dropped to 16 per year and continued to decline. He concluded that theory was a once useful heuristic that no longer serves a purpose in life history theory.
More recently, the panarchy theories of adaptive capacity and resilience promoted by C. S. Holling and Lance Gunderson have revived interest in the theory, and use it as a way of integrating social systems, economics, and ecology.
Writing in 2002, Reznick and colleagues reviewed the controversy regarding selection theory and concluded that:
Alternative approaches are now available both for studying life history evolution (e.g. Leslie matrix for an age-structured population) and for density-dependent selection (e.g. variable density lottery model).
| Biology and health sciences | Ecology | Biology |
1906880 | https://en.wikipedia.org/wiki/Rhizopus | Rhizopus | Rhizopus is a genus of common saprophytic fungi on plants and specialized parasites on animals. They are found in a wide variety of organic substances, including "mature fruits and vegetables", jellies, syrups, leather, bread, peanuts, and tobacco. They are multicellular. Some Rhizopus species are opportunistic human pathogens that often cause fatal disease called mucormycosis. This widespread genus includes at least eight species.
Rhizopus species grow as filamentous, branching hyphae that generally lack cross-walls (i.e., they are coenocytic). They reproduce by forming asexual and sexual spores. In asexual reproduction, spores are produced inside a spherical structure, the sporangium. Sporangia are supported by a large apophysate columella atop a long stalk, the sporangiophore. Sporangiophores arise among distinctive, root-like rhizoids. In sexual reproduction, a dark zygospore is produced at the point where two compatible mycelia fuse. Upon germination, a zygospore produces colonies that are genetically different from either parent.
Rhizopus oligosporus is used to make tempeh, a fermented food derived from soybeans.
Rhizopus oryzae is used in the production of alcoholic beverages in parts of Asia and Africa.
Rhizopus stolonifer (black bread mold) causes fruit rot on strawberry, tomato, and Sweet potato and is used in commercial production of fumaric acid and cortisone.
Various species, including R. stolonifer, may cause soft rot in sweet potatoes and Narcissus.
Rhizopus helps in nutrient development since this species is grown in soil it ferments the fruits and vegetable in the soil inhibiting the growth and develops certain pathogens that inhibits the growth of toxigenic fungus. In addition to that, there is even a type of Rhizopus (Rhizopus microsporus-fermented soybean tempe) that has proven to reduce colon carcinogenesis in rats by elevating factors of mucins, immunoglobulin A, and organic acids and give protection to piglets from Escherichia coli-infection by inhibiting adhesion to the intestinal membranes.
Phylogeny
The mating analysis has also been found which was comparative that this species structure is flexible in comparison with other species in the same genus. The topology length of the species genome is found to be three times bigger with the species.
Species
| Biology and health sciences | Basics | Plants |
2645795 | https://en.wikipedia.org/wiki/Taeniasis | Taeniasis | Taeniasis is an infection within the intestines by adult tapeworms belonging to the genus Taenia. There are generally no or only mild symptoms. Symptoms may occasionally include weight loss or abdominal pain. Segments of tapeworm may be seen in the stool. Complications of pork tapeworm may include cysticercosis.
Types of Taenia that cause infections in humans include Taenia solium (pork tapeworm), Taenia saginata (beef tapeworm), and Taenia asiatica (Asian tapeworm). Taenia saginata is due to eating contaminated undercooked beef while Taenia solium and Taenia asiatica is from contaminated undercooked pork. Diagnosis is by examination of stool samples.
Prevention is by properly cooking meat. Treatment is generally with praziquantel, though niclosamide may also be used. Together with cysticercosis, infections affect about 50 million people globally. The disease is most common in the developing world. In the United States fewer than 1,000 cases occur annually.
Signs and symptoms
Taeniasis generally has few or no symptoms. It takes about 8 weeks from infection for adult worms to form and can last for years without treatment.
Infection may be suspected when a portion of the worm is passed in the stool. It is not generally fatal.
Pork tapeworm
Infection in the intestines by the adult T. solium worms is normally asymptomatic. Heavy infection can result in anaemia and indigestion.
A complication, known as cysticercosis, may occur if the eggs of the pork tapeworm are eaten. This typically occurs from vegetables or water contaminated by feces from someone with pork tapeworm taeniasis. The eggs enter the intestine where they develop into larvae which then enter the bloodstream and invade host tissues. This is the most frequent and severe disease caused by any tapeworm. It can lead to headaches, dizziness, seizures, dementia, hypertension, lesions in the brain, blindness, tumor-like growths, and low eosinophil levels. It is a cause of major neurological problems, such as hydrocephalus, paraplegy, meningitis, and death.
Beef tapeworm
Taenia saginata infection is usually asymptomatic, but heavy infection causes weight loss, dizziness, abdominal pain, diarrhea, headaches, nausea, constipation, chronic indigestion, and loss of appetite. It can cause antigen reaction that induce allergic reaction. It is also a rare cause of ileus, pancreatitis, cholecystitis, and cholangitis.
Asian tapeworm
Taenia asiatica is also usually asymptomatic. It is unclear if T. asiatica can cause cysticercosis.
In pigs, the cysticercus can produce cysticercosis. Cysts develop in liver and lungs. (T. saginata does not cause cysticercosis.)
Transmission
Taeniasis is contracted after eating undercooked pork or beef that contains the larvae. The adult worms develop and live in the lumen of the intestine. They acquire nutrients from the intestine. The gravid proglottids, body segments containing fertilised eggs, are released in the faeces.
If consumed by an intermediate host such as a cow or pig, they hatch within the duodenum to become larvae, penetrate through the intestinal wall into nearby blood vessels, and enter the bloodstream. Once they reach organs such as the skeletal muscles, liver, or lungs, the larvae then develop into a cyst, a fluid-filled cysticercus. These contaminated tissues are then consumed through raw or undercooked meat.
Cysticercosis occurs when contaminated food, water, or soil that contains T. solium eggs is eaten.
Diagnosis
Diagnosis of taeniasis is mainly using stool sample, particularly by identifying the eggs. However, this has limitations at the species level because tapeworms have similar eggs. Examination of the scolex or the gravid proglottids can resolve the exact species. But body segments are not often available, therefore, laborious histological observation of the uterine branches and PCR detection of ribosomal 5.8S gene are sometimes necessary. Ziehl–Neelsen stain is also used for T. saginata and T. solium, in most cases only the former will stain, but the method is not entirely reliable. Loop-mediated isothermal amplification (LAMP) is highly sensitive (~2.5 times that of multiplex PCR), without false positives, for differentiating the taenid species from faecal samples.
To date the most relevant test for T. asiatica is by enzyme-linked immunoelectrotransfer blot (EITB). EITB can effectively identify asiatica from other taenid infections since the serological test indicates an immunoblot band of 21.5 kDa exhibited specifically by T. asiatica. Even though it gives 100% sensitivity, it has not been tested with human sera for cross-reactivity, and it may show a high false positive result.
Prevention
Prevention efforts include properly cooking meat, treating active cases in humans, vaccinating and treating pigs against the disease, stricter meat-inspection standards, health education, improved sanitation, and improved pig-raising practices.
Preventing human faeces from contaminating pig feeds also plays a role. Infection can be prevented with proper disposal of human faeces around pigs, cooking meat thoroughly and/or freezing the meat at −10 °C for 5 days. Contaminated hands are the primary method of transmission for human cysticercosis, especially in populations like food handlers.
Proper cooking of meat is an effective prevention. For example, cooking (56 °C for 5 minutes) of beef viscera destroys cysticerci. Refrigeration, freezing (−10 °C for 9 days) or long periods of salting is also lethal to cysticerci.
Treatment
Praziquantel is the treatment of choice. Usual treatments are with praziquantel (5–10 mg/kg, single-administration) or niclosamide (adults and children over 6 years: 2 g, single-administration after a light breakfast, followed after 2 hours by a laxative; children aged 2–6 years: 1 g; children under 2 years: 500 mg). One study showed albendazole is effective against animal beef tapeworm cysticercosis. Mepacrine is quite effective but has adverse effects in humans.
Epidemiology
The total global infection is estimated to be between 40 and 60 million people. In the US, the incidence of infection is low, but 25% of cattle sold are still infected.
Regions
Taeniasis is predominantly found in Asia, Africa, and Latin America, particularly on farms in which pigs are exposed to human excrement. At a low level though, it occurs everywhere where beef and pork are eaten, even in countries with strict sanitation policies such as the United States. Taenia saginata is relatively common in Africa, some parts of Eastern Europe, the Philippines, and Latin America. It is most prevalent in Sub-Saharan Africa and the Middle East. Taenia asiatica is restricted to East Asia, including Taiwan, Korea, Indonesia, Nepal, Thailand and China.
| Biology and health sciences | Helminthic diseases and infestations | Health |
12978962 | https://en.wikipedia.org/wiki/Plane%20table | Plane table | A plane table (plain table prior to 1830) is a device used in surveying, site mapping, exploration mapping, coastal navigation mapping, and related disciplines to provide a solid and level surface on which to make field drawings, charts and maps. The early use of the name plain table reflected its simplicity and plainness rather than its flatness.
"Plane" refers to the table being both flat levelled (horizontal).
History
The earliest mention of a plane table dates to 1551 in Abel Foullon's "Usage et description de l'holomètre", published in Paris. However, since Foullon's description was of a complete, fully developed instrument, it must have been invented earlier.
A brief description was also added to the 1591 edition of Digge's Pantometria. The first mention of the device in English was by Cyprian Lucar in 1590.
Some have credited Johann Richter, also known as Johannes Praetorius, a Nuremberg mathematician, in 1610 with the first plane table, but this appears to be incorrect.
The plane table became a popular instrument for surveying. Its use was widely taught. Some considered it a substandard instrument compared to other devices such as the theodolite, since it was relatively easy to use. By allowing the use of graphical methods rather than mathematical calculations, it could be used by those with less education than other instruments. The addition of a camera to the plane table, as was done from 1890 by Sebastian Finsterwalder in conjunction with a phototheodolite, established photogrammetry in spatial and temporal surveying.
Construction
A plane table consists of a smooth table surface mounted on a sturdy base. The connection between the table top and the base permits one to level the table precisely, using bubble levels, in a horizontal plane. The base, a tripod, is designed to support the table over a specific point on land. By adjusting the length of the legs, one can bring the table level regardless of the roughness of the terrain.
Usage
In use, a plane table is set over a point and brought to precise horizontal level. A drawing sheet is attached to the surface and an alidade is used to sight objects of interest. The alidade, in modern examples of the instrument a rule with a telescopic sight, can then be used to construct a line on the drawing that is in the direction of the object of interest.
By using the alidade as a surveying level, information on the topography of the site can be directly recorded on the drawing as elevations. Distances to the objects can be measured directly or by the use of stadia marks in the telescope of the alidade.
| Technology | Surveying tools | null |
12981660 | https://en.wikipedia.org/wiki/CMB%20cold%20spot | CMB cold spot | The CMB Cold Spot or WMAP Cold Spot is a region of the sky seen in microwaves that has been found to be unusually large and cold relative to the expected properties of the cosmic microwave background radiation (CMBR). The "Cold Spot" is approximately 70 μK (0.00007 K) colder than the average CMB temperature (approximately 2.7 K), whereas the root mean square of typical temperature variations is only 18 μK. At some points, the "cold spot" is 140 μK colder than the average CMB temperature.
The radius of the "cold spot" subtends about 5°; it is centered at the galactic coordinate , (equatorial: α = , δ = ). It is, therefore, in the Southern Celestial Hemisphere, in the direction of the constellation Eridanus.
Typically, the largest fluctuations of the primordial CMB temperature occur on angular scales of about 1°. Thus a cold region as large as the "cold spot" appears very unlikely, given generally accepted theoretical models. Various alternative explanations exist, including a so-called Eridanus Supervoid or Great Void that may exist between us and the primordial CMB (foreground voids can cause cold spots against the CMB). Such a void would affect the observed CMB via the integrated Sachs–Wolfe effect, and would be one of the largest structures in the observable universe. This would be an extremely large region of the universe, roughly 150 to 300 Mpc or 500 million to one billion light-years across and 6 to 10 billion light years away, at redshift , containing a density of matter much smaller than the average density at that redshift.
Discovery and significance
In the first year of data recorded by the Wilkinson Microwave Anisotropy Probe (WMAP), a region of sky in the constellation Eridanus was found to be colder than the surrounding area. Subsequently, using the data gathered by WMAP over 3 years, the statistical significance of such a large, cold region was estimated. The probability of finding a deviation at least as high in Gaussian simulations was found to be 1.85%. Thus it appears unlikely, but not impossible, that the cold spot was generated by the standard mechanism of quantum fluctuations during cosmological inflation, which in most inflationary models gives rise to Gaussian statistics. The cold spot may also, as suggested in the references above, be a signal of non-Gaussian primordial fluctuations.
Some authors called into question the statistical significance of this cold spot.
In 2013, the CMB Cold Spot was also observed by the Planck satellite at similar significance, discarding the possibility of being caused by a systematic error of the WMAP satellite.
Possible causes other than primordial temperature fluctuation
The large 'cold spot' forms part of what has been called an 'axis of evil' (so-called because it was unexpected to see a structure like this).
Supervoid
One possible explanation of the cold spot is a huge void between us and the primordial CMB. A region cooler than surrounding sightlines can be observed if a large void is present, as such a void would cause an increased cancellation between the "late-time" integrated Sachs–Wolfe effect and the "ordinary" Sachs–Wolfe effect. This effect would be much smaller if dark energy were not stretching the void as photons went through it.
Rudnick et al. found a dip in NVSS galaxy number counts in the direction of the Cold Spot, suggesting the presence of a large void. Since then, some additional works have cast doubt on the "supervoid" explanation. The correlation between the NVSS dip and the Cold Spot was found to be marginal using a more conservative statistical analysis. Also, a direct survey for galaxies in several one-degree-square fields within the Cold Spot found no evidence for a supervoid. However, the supervoid explanation has not been ruled out entirely; it remains intriguing, since supervoids do seem capable of affecting the CMB measurably.
A 2015 study shows the presence of a supervoid that has a diameter of 1.8 billion light years and is centered at 3 billion light-years from our galaxy in the direction of the Cold Spot, likely being associated with it. This would make it the largest void detected, and one of the largest structures known. Later measurements of the Sachs–Wolfe effect show too its likely existence.
Although large voids are known in the universe, a void would have to be exceptionally vast to explain the cold spot, perhaps 1,000 times larger in volume than expected typical voids. It would be 6 billion–10 billion light-years away and nearly one billion light-years across, and would be perhaps even more improbable to occur in the large-scale structure than the WMAP cold spot would be in the primordial CMB.
A 2017 study reported surveys showing no evidence that associated voids in the line of sight could have caused the CMB Cold Spot and concluded that it may instead have a primordial origin.
One important thing to confirm or rule out the late time integrated Sachs–Wolfe effect is the mass profile of galaxies in the area as ISW effect is affected by the galaxy bias which depends on the mass profiles and types of galaxies.
In December 2021, the Dark Energy Survey (DES), analyzing their data, put forward more evidence for the correlation between the Eridanus supervoid and the CMB cold spot.
Cosmic texture
In late 2007, (Cruz et al.) argued that the Cold Spot could be due to a cosmic texture, a remnant of a phase transition in the early Universe.
Parallel universe
A controversial claim by Laura Mersini-Houghton is that it could be the imprint of another universe beyond our own, caused by quantum entanglement between universes before they were separated by cosmic inflation. Laura Mersini-Houghton said, "Standard cosmology cannot explain such a giant cosmic hole" and made the hypothesis that the WMAP cold spot is "... the unmistakable imprint of another universe beyond the edge of our own." If true, this provides the first empirical evidence for a parallel universe (though theoretical models of parallel universes existed previously). It would also support string theory. The team claims that there are testable consequences for its theory. If the parallel-universe theory is true, there will be a similar void in the Celestial sphere's opposite hemisphere (which New Scientist reported to be in the Southern celestial hemisphere; the results of the New Mexico array study reported it as being in the Northern).
Other researchers have modeled the cold spot as potentially the result of cosmological bubble collisions, again before inflation.
A sophisticated computational analysis (using Kolmogorov complexity) has derived evidence for a north and a south cold spot in the satellite data: "...among the high randomness regions is the southern non-Gaussian anomaly, the Cold Spot, with a stratification expected for the voids. Existence of its counterpart, a Northern Cold Spot with almost identical randomness properties among other low-temperature regions is revealed."
These predictions and others were made prior to the measurements (see Laura Mersini). However, apart from the Southern Cold Spot, the varied statistical methods in general fail to confirm each other regarding a Northern Cold Spot. The 'K-map' used to detect the Northern Cold Spot was noted to have twice the measure of randomness measured in the standard model. The difference is speculated to be caused by the randomness introduced by voids (unaccounted-for voids were speculated to be the reason for the increased randomness above the standard model).
Sensitivity to finding method
The cold spot is mainly anomalous because it stands out compared to the relatively hot ring around it; it is not unusual if one only considers the size and coldness of the spot itself. More technically, its detection and significance depends on using a compensated filter like a Mexican hat wavelet to find it.
| Physical sciences | Notable patches of universe | Astronomy |
22450812 | https://en.wikipedia.org/wiki/Bedform | Bedform | A bedform is a geological feature that develops at the interface of fluid and a moveable bed, the result of bed material being moved by fluid flow. Examples include ripples and dunes on the bed of a river. Bedforms are often preserved in the rock record as a result of being present in a depositional setting. Bedforms are often characteristic to the flow parameters, and may be used to infer flow depth and velocity, and therefore the Froude number.
Bedforms Initiation
Bedforms are omnipresent in many environments (e.g., fluvial, eolian, glaciofluvial, deltaic and deep sea), although there is still some debate on how they develop. There are two separate, though not mutually exclusive, models of bedform initiation: defect initiation and instantaneous initiation.
Defect Initiation
The defect theory proposes that the turbulent sweeps that are generated in turbulent flows entrain sediment that upon deposition generates defects in a non-cohesive material. These deposits then propagate downstream via a flow separation process, thus developing bedform fields. The origin of the defects is thought to be linked to packets of hairpin vortex structures. These coherent turbulent structures give rise to entrainment corridors on the mobile bed, forming grain lineations that interact with the low-speed streaks generating an agglomeration of grains. Once a critical height of grains is reached, flow separation occurs over the new structure. Sediment will be eroded close from the reattachment point and deposited downstream creating a new defect. This new defect will thus induce formation of another defect and the process will continue, propagating downstream while the accumulations of grains quickly evolve into small bedforms.
Instantaneous Initiation
In general, the defect propagation theory plays a bigger role at low sediment transport rates since for high rates defects maybe washed away and bedforms generally initiated across the entire bed spontaneously. Venditti et al. (2005) report that instantaneous initiation begins with the formation of a cross-hatch pattern, which leads to chevron-shaped forms that migrate independently of the pattern structure. This chevron-like structure reorganizes to form the future crest lines of the bedforms. Venditti et al. (2006), based on the earlier model by Liu (1957), proposed that instantaneous initiation is a manifestation of an interfacial hydrodynamic instability of Kelvin-Helmholtz type between a highly active pseudofluid sediment layer and the fluid above it. In addition, Venditti et al. (2005) imply that there is no linkage between the instantaneous initiation and coherent turbulent flow structures, since spatially- and temporally-random events should lock in place to generate the cross-hatch pattern. Moreover, there is no clear explanation of the effect of turbulence in the formation of bedforms since bedforms may also occur under laminar flows . It is important to note, that laminar-generated bedform studies used the temporally-averaged flow conditions to determine the degree of turbulence, indicating Reynolds number in the laminar regime. However, instantaneous process, such as burst and sweeps, which are infrequent at low Reynolds number but still present, can be the driving mechanisms to generate the bedforms. The generation of bedforms in laminar flows is still a topic of debate within the scientific community, since if true, it suggests that there should be other processes for defect development other than the one suggested by Best (1992). This alternative model for bedform development at low sediment transport rates should explain the generation of defects and bedforms for cases where the flow is not turbulent.
Bedform phase diagrams
Phase or stability diagrams are defined as graphs that show the regimes of existence of one or more stable bed states. The stability of the bed can be defined when the bedform is in equilibrium and does not change in time for the same flow condition. This invariance over time must not be confused with a static morphology or frozen equilibrium; on the contrary, the bed moves and adjusts in a dynamic equilibrium with the flow and sediment transport for that particular condition. These phase diagrams are used for two main purposes: i) for prediction of bed states in a known flow and sediment transport condition, and, ii) as a tool for the reconstruction of paleoenvironments from a known bed state or sedimentary structure. Despite the great utility of such diagrams, they are very difficult to construct, making them either incomplete or very hard to interpret. This complexity lies in the number of variables needed to quantify the system.
Bedforms vs. flow
Typical unidirectional bedforms represent a specific flow velocity, assuming typical sediments (sands and silts) and water depths, and a chart such as below can be used for interpreting depositional environments, with increasing water velocity going down the chart.
This chart is for general use, because changes in grain size and flow depth can change the bedform present and skip bedforms in certain scenarios. Bidirectional environments (e.g. tidal flats) produce similar bedforms, but the reworking the sediments and opposite directions of flow complicates the structures.
This bed form sequence can also be illustrated diagrammatically:
Types of Bedforms
Lower Plane Bed
"Lower plane bed" refers to the flat configuration the bed of a river that is produced in via low rates of sediment transport.
Upper Plane Bed
"Upper plane bed" features are flat and characterized by a unidirectional flow with high rates of sediment transport as both bed load and suspended load. Upper plane bed conditions can produce parting current lineations, which are typically subtle streaks on the bed surface due to the high energy flow.
| Physical sciences | Sedimentology | Earth science |
4857087 | https://en.wikipedia.org/wiki/Aluminosilicate | Aluminosilicate | Aluminosilicate refers to materials containing anionic Si-O-Al linkages. Commonly, the associate cations are sodium (Na+), potassium (K+) and protons (H+). Such materials occur as minerals, coal combustion products and as synthetic materials, often in the form of zeolites. Both synthetic and natural aluminosilicates are of technical significance as structural materials, catalysts, and reagents.
Important representatives
Feldspar is a common tectosilicate aluminosilicate mineral made of potassium, sodium, and calcium cations surrounded by a negatively charged network of silicon, aluminium and oxygen atoms.
Many aluminosilicates are synthesized by reactions of silicates, aluminates, and other compounds. They have the general formula where M+ is usually H+ and Na+. The Si/Al ratio is variable, which provides a means to tune the properties. Many of these materials are porous and exhibit properties of industrial value. Naturally occurring microporous, hydrous aluminosilicate minerals are also referred to as zeolites.
| Physical sciences | Silicate minerals | Earth science |
4857714 | https://en.wikipedia.org/wiki/Cumacea | Cumacea | Cumacea is an order of small marine crustaceans of the superorder Peracarida, occasionally called hooded shrimp or comma shrimp. Their unique appearance and uniform body plan makes them easy to distinguish from other crustaceans. They live in soft-bottoms such as mud and sand, mostly in the marine environment. There are more than 1,500 species of cumaceans formally described. The species diversity of Cumacea increases with depth.
Anatomy
Cumaceans have a strongly enlarged cephalothorax with a carapace, a slim abdomen, and a forked tail. The length of most species varies from .
The carapace of a typical cumacean is composed of several fused dorsal head parts and the first three somites of the thorax. This carapace encloses the appendages that serve for respiration and feeding. In most species, there are two eyes at the front side of the head shield, often merged into a single dorsal eye lobe. The five posterior somites of the thorax form the pereon. The pleon (abdomen) consists of six cylindrical somites.
The first antenna (antennule) has two flagella, the outer flagellum usually being longer than the inner one. The second antenna is strongly reduced in females, and consists of numerous segments in males.
Cumaceans have six pairs of mouthparts: one pair of mandibles, one pair of maxillules, one pair of maxillae and three pairs of maxillipeds.
Ecology
Cumaceans are mainly marine crustaceans. However, some species can survive in water with a lower salinity, like brackish water (e.g. estuaries). In the Caspian Sea they even reach some rivers that flow into it. A few species live in the intertidal zone.
Most species live only one year or less, and reproduce twice in their lifetime. Deep-sea species have a slower metabolism and presumably live much longer.
Cumaceans feed mainly on microorganisms and organic material from the sediment. Species that live in the mud filter their food, while species that live in sand browse individual grains of sand. In the genus Campylaspis and a few related genera, the mandibles are transformed into piercing organs, which can be used for predation on foraminiferans and small crustaceans.
Many shallow-water species show a diurnal cycle, with males emerging from the sediment at night and swarming to the surface.
Importance
Like Amphipoda, cumaceans are an important food source for many fishes. Therefore, they are an important part of the marine food chain. They can be found on all continents.
Reproduction and development
Cumaceans are a clear example of sexual dimorphism: males and females differ significantly in their appearance. Both sexes have different ornaments (setation, knobs, and ridges) on their carapace. Other differences are the length of the second antenna, the existence of pleopods in males, and the development of a marsupium (brood pouch) in females. There are generally more females than males, and females are also larger than their male counterparts.
Cumaceans are epimorphic, which means that the number of body segments does not change during development. This is a form of incomplete metamorphosis. Females carry the embryos in their marsupium for some time. The larvae leave the marsupium in the manca stage, in which they are almost fully grown and are only missing their last pair of pereiopods.
History of research
The order Cumacea has been known since 1780, when Ivan Ivanovich Lepechin described the species "Oniscus scorpioides" (now Diastylis scorpioides). At the time, many scientists thought that the cumaceans were larval stages of decapods. In 1846, they were recognised as a separate order by Henrik Nikolaj Krøyer. Twenty-five years later, about fifty different species had been described, and currently there are more than 1,500 described species. The German zoologist Carl Wilhelm Erich Zimmer studied the order Cumacea very intensively.
Fossil record
The fossil record of cumaceans is very sparse, but extends back into the Mississippian age. Fossil Cumaceans from the early Jurassic scarcely differ from living forms (Bacescu & Petrescu 1999).
Eobodotria muisca was found in 2019 in strata from the Middle Cretaceous of Colombia. Exceptional details such as the gut, mouth parts, pereopods, setae bearing uropods, antenna with developed flagella, and even small eyes with ommatidia were preserved. Eobodotria straddles a gap of almost 165 million years in the fossil record of sea commas, providing a reliable calibration point for phylogenetic studies. This species is considered the first certain representative of crown Cumacea.
Taxonomy
Cumaceans belong to the superorder Peracarida, within the class Malacostraca. The order Cumacea is subdivided into 8 families, 141 genera, and 1,523 species:
Bodotriidae Scott, 1901 (379 species in 36 genera)
Ceratocumatidae Calman, 1905 (10 species in 2 genera)
Diastylidae Bate, 1856 (318 species in 22 genera)
Gynodiastylidae Stebbing, 1912 (106 species in 12 genera)
Lampropidae Sars, 1878 (114 species in 15 genera)
Leuconidae Sars, 1878 (139 species in 16 genera)
Nannastacidae Bate, 1866 (426 species in 25 genera)
Pseudocumatidae Sars, 1878 (30 species in 12 genera)
One species is also placed incertae sedis in the order.
| Biology and health sciences | Malacostraca | Animals |
4863813 | https://en.wikipedia.org/wiki/Halkieriid | Halkieriid | The halkieriids are a group of fossil organisms from the Lower to Middle Cambrian. Their eponymous genus is Halkieria , which has been found on almost every continent in Lower to Mid Cambrian deposits, forming a large component of the small shelly fossil assemblages. The best known species is Halkieria evangelista, from the North Greenland Sirius Passet Lagerstätte, in which complete specimens were collected on an expedition in 1989. The fossils were described by Simon Conway Morris and John Peel in a short paper in 1990 in the journal Nature. Later a more thorough description was undertaken in 1995 in the journal Philosophical Transactions of the Royal Society of London and wider evolutionary implications were posed.
The group is sometimes equated to Sachitida, although as originally envisaged, this group includes the wiwaxiids and is thus equivalent to the Halwaxiida.
History of discovery
Armor plates called "sclerites" had long been known as elements of the small shelly fossils, and detailed analysis showed that some of these belonged to the same animal and how they fitted together. The first articulated specimens of Halkieria evangelista, with all their hard parts together, were collected in 1989 from the Sirius Passet lagerstätte in Greenland, and were described in 1990 by Simon Conway Morris and John S. Peel. H. evangelista is used as a model for identifying and reconstructing as halkieriids other similar shells and sclerites; its epithet evangelista reflects its power to explain the Lower Cambrian fossil record.
Description of the fossils
Features shared by Halkieria and Australohalkieria
Only armor-like sclerites of Australohalkieria have been found, and much of the analysis assumes that these animals were similar to Halkieria. However the sclerites are so similar that this assumption looks fairly safe. In both genera the sclerites are of the type called "coelosclerites", which have a mineralized shell around a space originally filled with organic tissue, and which show no evidence of growth by adding material round the outside. Both genera also have sclerites of three different shapes: "palmate", flat and shaped rather like a maple leaf, which are generally the smallest; "cultrate", flat but shaped like knife blades; and "siculate", which are about the same size as the cultrates but are spine-shaped and like rather squashed cylinders. In both Halkieria and Australohalkieria the palmate and cultrate sclerites have prominent ribs, and are fairly flat except for 90° bends at the bases, which indicate that they fitted snugly against the animals' bodies. The siculates mostly lack ribs and appear to have projected away from the body at angles between about 45° and 90°.
Halkieria evangelista
The animals looked like slugs in chain mail - to long, bilaterally symmetric, flattened from top to bottom and unarmored on the bottom. Very near each end there is a shell plate with prominent growth lines rather like the growth rings of trees. The rest of the upper surface was covered with about 2,000 sclerites that overlapped each other like tiles and formed three zones with sclerites of different shapes: "palmates", shaped rather like maple leaves, ran along the center of the back between the shell plates; blade-shaped "cultrates" lay on either side of the palmates and pointing towards the middle of the upper surface; and slim, sickle-shaped "siculates" covered the outer edges. The sclerites bore a wide central cavity, and (at least in some specimens) finer lateral canals. As the animals grew, the shell plates grew by adding material to the outer edges. Individual sclerites stayed the same size; since the cultrate sclerites form a pattern that is constant in all fairly complete specimens, the old ones that were too small may have been shed and replaced by larger ones as the animals grew. The sclerites seem to have grown by basal secretion. There are traces of thin ribs between the sclerites and the skin.
The shellplates and the sclerites were probably made of calcium carbonate originally; it has been suggested on the basis of how they were preserved that they may have been wholly organic, but this is less likely since fossils of non-calcified organisms are usually thin films while Halkeieria fossils are three-dimensional like those of trilobites and hyoliths - in fact several specimens show curvature in the horizontal plane, which suggests that the muscles associated with the sclerites were still present at the
time of burial
The sole was soft and probably muscular. Since Halkieria was unsuited to swimming and had no obvious adaptations for burrowing, it must have lived on the sea-floor, "walking" by making its muscular sole ripple. The backward-projecting siculate sclerites may have improved its grip by preventing it from slipping backwards. Some specimens have been found partially rolled up, rather like a pillbug, and in this position the cultrate sclerites projected outwards, which probably deterred predators. It is difficult to determine the functions of the cap-shaped shells at either end of the animal, as the sclerites appear to have offered adequate protection. Scars on the inner surface of the front shell may indicate that it provided an attachment for internal organs. In one specimen the rear shell appears to have rotated by about 45° before fossilization, which suggests there was a cavity underneath, which may have housed gills.
Traces of a gut have been found in the rear halves of some fossils. Parts of one specimen have been interpreted as a radula, the toothed chitinous tongue that is the signature feature of molluscs, but in this specimen the edge of the "scleritome", i.e. coat of sclerites, is folded and the putative radula could be a group of dislocated siculate sclerites.
Australohalkieria superstes
The name of the most complete and abundant Australian find means "Southern Halkieria the Survivor" because it proves that halkieriids survived the end-Botomian extinction. The sclerites assigned to this species are convex on the upper surface and concave on the lower. They may also curve within their own plane, and they overlap so that the concave side of each is partly covered by the convex side of the next one. The internal cavity within Australohalkieria is more complicated than the simple tube in Halkieria; about half-way up the sclerite, the cylindrical tube splits into a pair of longitudinal canals, with the central canal flattening; the canals don't seem to be connected. The walls also have a different microscopic structure.
In A. superstes the central canals of sclerites are flattened on their upper surfaces, and this produces a depression on the upper surface of the tip. The surface of this depression is not mineralized, which suggests the depression may have helped the animals' sense of smell by letting chemicals in the water penetrate the exposed skin. The phosphatic coating on sclerites of A. superstes has features that suggest they were originally covered by a thin organic skin. An outer organic layer has also been found on sclerites of the chancelloriids, sessile organisms that are thought to have looked rather like cacti. If halkieriids were early molluscs, the outer layers of the sclerites may have been similar to the periostracum of some modern molluscs.
The sclerites of A. superstes have right- and left-handed variants which are equally abundant, which suggests that A. superstes was bilaterally symmetrical. All of the sclerites were tiny: the palmate ones ranged from to in length, and the cultrates from to . The siculates fall into two groups: those with a shallow S-curve at the base, which range from to in length, and often have a slight twist at the base; and those with a 45° and 90° bend at the base and are to long.
Scleritomes of Early Cambrian halkieriids have many more palmate and cultrate than siculate sclerites. On the other hand, siculate sclerites of A. superstes are more abundant than either cultrate or palmate sclerites; in fact palmate sclerites are rare. Possibly some process after death removed many of the palmates and some of the cultrates, but it is more likely that in A. superstes the part of the scleritome, or "coat of mail", closest to the sea-bed was larger relative to the lateral and dorsal zones further up and towards the center. A. superstes sclerites are also about one-third the size of Early Cambrian halkieriid sclerites. Since the Georgina assemblage includes larger fossils and most Early Cambrian halkieriids are preserved by the same method, phosphatization, it is unlikely that preservational bias has produced an unrepresentative sample. Possible explanations for the small size of A. superstes sclerites include: the individual(s) represented in the Georgina assemblage were juveniles; their scleritomes were composed of many more sclerites than those of Early Cambrian halkieriids; or the species itself was relatively small.
No shells that might be assigned to halkieriids have been found in the Georgina Basin. This does not prove that Australohalkieria lacked shells, as shells of Halkieria are rarely found.
Australohalkieria parva
This species, whose name means "Small Southern Halkieria", was first described in 1990. Like A. superstes, its sclerites have undivided longitudinal canals and a very similar structure to their walls wall, but A. parva has sclerites whose central canals are not flattened.
Other halkieriid fossils from Australia
The other sclerites from the Georgina Basin are different enough to be excluded from Australohalkieria superstes, but are not sufficiently abundant to provide enough detail for them to be classified. One type is very similar to those of A.superstes, even having a two-pronged tip, but the middle canal is not flattened. The other has a flattened central canal and no longitudinal canals, and may represent an additional Middle Cambrian halkieriid genus, distinct from Australohalkieria and from the Early Cambrian Halkieria.
Siphogonuchitids
Siphogonuchitids have two sclerite morphs as well as their shell(s), thus may have had a simpler scleritome than Halkieria and its ilk, concordant with the sclerites' simpler internal anatomy.
The genera Siphogonuchites, Dabashanites, Lopochites, and Maikhanella all seem to represent components of the Siphogonuchites animal.
Sclerites of Drepanochites can be distinguished based on their aspect ratio.
Maikhanella is shell formed of Siphogonuchites sclerites that are fused together with a calcified matrix. Juvenile shells appear not to incorporate sclerites.
The central cavity of the Siphogonuchites sclerite is simple, with no lateral chambers attached.
Ninellids
The ninellids, typified by Ninella, are a Lower Cambrian group that had an even simpler scleritome, with only one sclerite type (although variation in the morphology of the sclerites is observed, and left- and right-sided sclerites exist). Their sclerites are hooked or scoop-like, and are very similar to halkieriid or siphonogunuchitid sclerites; they were hollow and calcareous and had a ridged upper surface.
Hippopharangites
Hippopharangites has sclerites with a broad central cavity and small pores opening through the shell wall, equivalent to the lateral chambers of other halkieriids (and the aesthete canals of Chitons?)
This genus is the closest in form to Chancelloriid sclerites, and is thus used to support the union of halkieriids and chancelloriids as Coeloscleritophora.
Lomasulcachites
Lomasulcachites is a further genus known from sclerites alone.
Sachites
Sachites Meshkova 1969 is a genus that comprised spiny sclerites; many Sachites specimens are now referred to other halkieriid taxa.
Although believed to be related to the halkieriids, a chancelloriid affinity has more recently been proposed.
Sinosachites
Sinosachites is a genus of 'halkieriid' known only from sclerites; these have internal chambers that are sub-perpendicular to the central canal, to which they are connected by narrow channels. The chambers are the same diameter, ~40 μm, as the longitudinal canals in Australohalkieria; their greater number and arrangement as lateral rather than longitudinal bodies reflects the greater size of the Sinosachites sclerites, which measure about 1–2 mm in length.
The sclerites are synonymous with Thambetolepis, which was originally described from Australia.
Left-hand and right-hand sclerites exist, so the animal was bilaterally symmetrical; as in Halkieria, palmate, cultrate and siculate sclerite morphologies exist.
Oikozetetes
Oikozetetes is known only from two types of cap-shaped shell found in the Burgess Shale and dated to about . The two types are thought to be the front and rear shells of a halkieriid.
They were probably calcareous while the organism was alive (although diagenesis sometimes replaces the original mineral with another, such as silica). It is thought to also have borne an armour coat consisting of biomineralised sclerites, like Halkieria. These are never found in direct association with the shells, but there are many biostratinomic processes which could account for this fact.
The lower Cambrian taxon Ocruranus (=Eohalobia) is putatively equivalent to the shells of Oikozetetes and seemingly belonged to a halkieriid-type body, although an intermediate valve suggests a Palaeoloricate-like body form.
Occurrence
The only reasonably complete specimens, of Halkieria evangelista, were found in the Sirius Passet lagerstätte in Greenland. Fragments which are confidently classified as belonging to halkieriids have been found in China's Xinjiang province and Australia's Georgina Basin, while shells of a possible halkieriid have been found in Canada's Burgess Shale. Halkieriid-like armor plates, called "sclerites" have been found in many other places as part of the small shelly fauna.
The earliest known occurrences of Halkieriids sclerites, classified as Halkieria longa, date from the Purella antiqua Zone of the Upper Nemakit-Daldynian Stage in Siberia. The mass extinction at the end of the Cambrian period's Botomian age was thought to have wiped out most of the small shellies, including the halkieriids, but in 2004 Halkieriid fossils classified as Australohalkieria were reported from Mid-Cambrian rocks of the Georgina Basin in Australia. It is not known why this clade would have survived while other halkieriid clades apparently died. It may be significant that the only archaeocyathans known to have survived the end-Botomian extinction also occur in Gondwana, the old super-continent that embraced South America, Africa, India, Australia and Antarctica.
Halkieriids and other small shelly fossils are typically, although not always, preserved in phosphate, which may or may not have been their original mineral composition. Preservation by a covering of phosphate only seems to have been common during the early Cambrian, becoming rarer with time as a result of increased disturbance of sea-floors by burrowing animals. Hence it is possible that halkieriids and other small shelly fossils were alive earlier than the earliest known fossils and later than the latest known fossils — paleontologists call this kind of uncertainty the Signor–Lipps effect.
Species
Nearly all members of the genera Halkieria are known only from finds of isolated scaly sclerites:
Halkieria alata Duan, 1984
Halkieria amorpha Meshkova,1974
Halkieria bisulcata Qian et Yin, 1984
Halkieria costulata Meshkova, 1974
Halkieria curvativa Mambetov in Missarzhevsky and Mambetov, 1981
Halkieria deplanatiformis Mambetov in Missarzhevsky and Mambetov, 1981
Halkieria desquamata Duan, 1984
Halkieria directa Mostler, 1980
Halkieria elonga Qian et Yin, 1984
Halkieria equilateralis Qian et Yin, 1984
Halkieria folliformis Duan, 1984
Halkieria fordi Landing, 1991
Halkieria hexagona Mostler, 1980
Halkieria lata Mostler, 1980
Halkieria longa Qian, 1977
Halkieria longispinosa Mostler, 1980
Halkieria maidipingensis Qian, 1977
Halkieria mina Qian, Chen et Chen, 1979
Halkieria mira Qian et Xiao, 1984
Halkieria obliqua Poulsen, 1967
Halkieria operculus Qian, 1984
?Halkieria pennata He, 1981 [=?Halkieria sthenobasis Jiang in Luo et al., 1982]
Halkieria phylloidea He, 1981
Halkieria praeinguis Jiang in Luo et al., 1982
Halkieria projecta Bokova, 1985
Halkieria sacciformis Meshkova, 1969
Halkieria solida Mostler, 1980
Halkieria sthenobasis Jiang in Luo et al., 1982
Halkieria stonei Landing, 1989
Halkieria symmetrica Poulsen, 1967
Halkieria terastios Qian, Chen et Chen, 1979
Halkieria uncostata Qian et Yin, 1984
Halkieria undulata Wang, 1994
Halkieria ventricosa Mostler, 1980
Halkieria wangi Demidenko, 2010
Halkieria zapfei Mostler, 1980
At present, the structure of complete scleritome is known only for the single species named Halkieria evangelista from the Lower Cambrian of Greenland (Sirius Passet Formation).
Phylogenetic position of Halkieria
The evolutionary relationships of the halkieriids are a complex topic which is still being debated. Most of this debate is about their relationship to Wiwaxia and to the three major lophotrochozoan phyla — molluscs, annelids and brachiopods. The question of their relationship to an apparently much more primitive Cambrian group, the chancelloriids is also significant and may raise some difficult questions.
Relationship to Molluscs, Annelids and Brachiopods
In 1995 Conway Morris and Peel presented a cladogram based both on the fossils' features and on early 1990s research in molecular phylogeny, which is the application of cladistic analysis to DNA and RNA:
The siphogonotuchids, a group found in Earliest Cambrian rocks, were the "sister" group to all the rest. These are known only from isolated fragments.
The earliest halkieriids were a "sister" group to the molluscs, in other words descendants of a fairly closely related common ancestor. This relationship, they said, was supported by the muscular foot that most researchers assumed halkieriids had.
Another halkieriid genus, Thambetolepis / Sinosachites, was a "great aunt" of annelids and Wiwaxia was an "aunt" of annelids. Their claim of a close relationship between halkieriids and Wiwaxia was based on both groups' having sclerites divided into three concentric zones. The close relationship of Wiwaxia to annelids was based on the similarities Butterfield (1990) found between Wiwaxia sclerites and the bristles of polychaete annelids. Canadia is a Burgess Shale fossil that is widely agreed to be a polychaete.
Halkieria evangelista, which Conway Morris had found in Greenland's Sirius Passet lagerstätte, was a "sister" group" to brachiopods, animals whose modern forms have bivalve shells but differ from molluscs in having muscular stalks and a distinctive feeding apparatus, the lophophore. Brachiopods have bristles that are similar to those of annelids and hence to Wiwaxia sclerites, and hence to halkieriid sclerites. A brachiopod affinity seemed plausible because brachiopods pass through a larval phase that resembles a halkieriid, and some isolated fossil shells thought to belong to halkieriids had a brachiopod-like microstructure.
In 2003 Cohen, Holmer and Luter supported the halkieriid-brachiopod relationship, suggesting that brachiopods may have arisen from a halkieriid lineage that developed a shorter body and larger shells, and then folded itself and finally grew a stalk out of what used to be the back.
Vinther and Nielsen (2005) proposed instead that Halkieria was a crown group mollusc, in other words more similar to modern molluscs that to annelids, brachiopods or any intermediate groups. They argued that: Halkieria sclerites resembled those of the modern solenogaster aplacophoran shell-less molluscs (see ), of some modern polyplacophoran molluscs, which have several shell plates, and of the Ordovician polyplacophoran Echinochiton; Halkieria shells are more similar to the shells of conchiferan molluscs, since shells of both of these groups show no trace of the canals and pores seen in polyplacophoran shell plates; the bristles of brachiopods and annelids are similar to each other but not to Halkieria sclerites.
Caron, Scheltema, Schander and Rudkin (2006) also interpreted Halkieria as a crown group mollusc, with Wiwaxia and Odontogriphus as stem group molluscs, in other words "sister" and "aunt" of the crown group molluscs. Their main reason for regarding Halkieria as crown group molluscs is that both possessed armor mineralized with calcium carbonate. They treated Wiwaxia and Odontogriphus as stem group molluscs because in their opinion both possessed the distinctive molluscan radula, a chitinous toothed "tongue".
Also in 2006, Conway Morris criticized Vinther and Nielsen's (2005) classification of Halkieria as a crown group mollusc, on the grounds that the growth of the spicules in the aplacophorans and polyplacophorans is not similar to the method of growth deduced for the complex halkieriid sclerites; in particular, he said, the hollow spines of various molluscs are not at all like the halkieriid sclerites with their complex internal channels. Conway Morris repeated his earlier conclusion that halkieriids were close to the ancestors of both molluscs and brachiopods.
Butterfield (2006) accepted that Wiwaxia and Odontogriphus were closely related, but argued that they were stem-group polychaetes rather than stem-group molluscs. In his opinion the feeding apparatus of these organisms, which consisted of two or at most four rows of teeth, could not perform the functions of the "belt-like" molluscan radula with their numerous tooth-rows; the different tooth-rows in both Wiwaxia and Odontogriphus tooth-rows also have noticeably different shapes, while those of molluscan radulae are produced one after the other by the same group of "factory" cells and therefore are almost identical. He also regarded lines running across the middle region of Odontogriphus fossils as evidence of external segmentation, since the lines are evenly spaced and run exactly at right angles to the long axis of the body. As in his earlier papers, Butterfield emphasized the similarities of internal structure between Wiwaxia sclerites and the bristles of polychaetes, and the fact that polychaetes are the only modern organisms in which some of the bristles form a covering over the back.
Conway Morris and Caron (2007) published the first description of Orthrozanclus reburrus. This resembled the halkieriids in having concentric bands of sclerites, although only two and not mineralized; and one shell at what was presumed to be the front and which was similar in shape to Halkieria front shell. It also had long spines rather like those of Wiwaxia. Conway Morris and Caron regarded this creature as evidence that the "halwaxiids" were a valid taxon and were monophyletic, in other words shared a common ancestor with each other and with no other organism. They published two cladograms, representing alternative hypotheses about the evolution of the lophotrochozoa, the lineage that includes molluscs, annelids and brachiopods:
This is the more likely, although it falls apart if the organisms' characteristics are changed even slightly:
Kimberella and Odontogriphus are early, primitive molluscs, without sclerites or any kind of mineralized armor.
Wiwaxia, the siphogonotuchids, Orthrozanclus and Halkieria from a side-branch of the mollusc family tree, which diverged in that order. This would mean that: Wiwaxia was the first of them to have sclerites, which were unmineralized; the siphogonotuchids were the first to have mineralized sclerites, although the scleritome was simpler; halkieriids then develop more complex scleritomes, while in Orthrozanclus the scleritome became unmineralized again and the rear shell vanished or became so small that it has not been seen in fossils. This hypothesis faces the difficulty that siphogonotuchids appear in earlier rocks and have simpler scleritomes than the other three groups.
The annelids and brachiopods evolved from the other main branch of the family tree, which did not include the molluscs.
The alternative view is:
Kimberella and Odontogriphus are early, primitive lophotrochozoans.
The siphogonotuchids, Halkieria, Orthrozanclus and Wiwaxia form a group that is closer to the shared ancestor of annelids and brachiopods than it is to the molluscs. The siphogonotuchids are the first of the group to become distinctive, with two types of mineralized sclerites and a "shell" made of fused sclerites. Halkieriids had three types of sclerites and two one-piece shells. In Orthrozanclus the sclerites became unmineralized and in Wiwaxia the shells were lost.
The network of internal cavities within sclerites of the halkieriid Sinosachites have been likened to the aesthete canals in polyplacophora, strengthening the case for a molluscan affinity.
Relationship to chancelloriids
Porter (2008) revived an early 1980s idea that the sclerites of Halkieria are extremely similar to those of chancelloriids. These were sessile, bag-like, radially symmetric organisms with an opening at the top.
Since their fossils show no signs of a gut or other organs, they were originally classified as some kind of sponge. Butterfield and Nicholas (1996) argued that they were closely related to sponges on the grounds that the detailed structure of chancellorid sclerites is similar to that of fibers of spongin, a collagen protein, in modern keratose (horny) demosponges. However Janussen, Steiner and Zhu (2002) opposed this view, arguing that: spongin does not appear in all Porifera, but may be a defining feature of the demosponges; the silica-based spines of demosponges are solid, while chancellorid sclerites are hollow and filled with soft tissues connected to the rest of the animal at the bases of the sclerites; chancellorid sclerites were probably made of aragonite, which is not found in demosponges; sponges have loosely bound-together skins called pinacoderms, which are only one cell thick, while the skins of chancellorids were much thicker and shows signs of connective structures called belt desmosomes. In their opinion the presence of belt desmosomes made chancellorids members of the Epitheliazoa, the next higher taxon above the Porifera, to which sponges belong. They thought it was difficult to say whether chancellorids were members of the Eumetazoa, "true animals" whose tissues are organized into Germ layers: chancellorids' lack of internal organs would seem to exclude them from the Eumetazoa; but possibly chancellorids descended from Eumetazoans that lost these features after becoming sessile filter-feeders. There are intriguing hints that the Ediacaran genus Ausia may represent a halkieriid ancestor with strong similarity to the chancelloriids.
The coelosclerites ("hollow sclerites") of halkieriids and chancelloriids resemble each other at all levels: both have an internal "pulp cavity" and a thin external organic layer; the walls are made of the same material, aragonite; the arrangement of the aragonite fibers is in each is the same, running mainly from base to tip but with each being closer to the surface at the end nearest the tip. It is extremely improbable that totally unrelated organisms could have developed such similar sclerites independently, but the huge difference in the structures of their bodies makes it hard to see how they could be closely related. This dilemma may be resolved in various ways:
One possibility is that chancelloriids evolved from bilaterian ancestors but then adopted a sessile lifestyle and rapidly lost all unnecessary features. However the gut and other internal organs have not been lost in other bilaterians that lost their external bilateral symmetry, such as echinoderms, priapulids, and kinorhynchs.
On the other hand, perhaps chancelloriids are similar to the organisms from which bilaterians evolved. That would imply that the earliest bilaterians had similar coelosclerites. However, there are no fossils of such sclerites before , while Kimberella from was almost certainly a bilaterian, but shows no evidence of sclerites.
One solution to this dilemma may be that preservation of small shelly fossils by coatings of phosphate was common only for a relatively short time, during the Early Cambrian, and that coelosclerite-bearing organisms were alive several million years before and after the time of phosphatic preservation. In fact there are over 25 cases of phosphatic preservation between and , but only one between and .
Alternatively, perhaps the common ancestor of both chancelloriids and halkieriids had very similar but unmineralized coelosclerites, and some intermediate groups independently incorporated aragonite into these very similar structures.
| Biology and health sciences | Mollusks | Animals |
27147535 | https://en.wikipedia.org/wiki/Transport%20phenomena | Transport phenomena | In engineering, physics, and chemistry, the study of transport phenomena concerns the exchange of mass, energy, charge, momentum and angular momentum between observed and studied systems. While it draws from fields as diverse as continuum mechanics and thermodynamics, it places a heavy emphasis on the commonalities between the topics covered. Mass, momentum, and heat transport all share a very similar mathematical framework, and the parallels between them are exploited in the study of transport phenomena to draw deep mathematical connections that often provide very useful tools in the analysis of one field that are directly derived from the others.
The fundamental analysis in all three subfields of mass, heat, and momentum transfer are often grounded in the simple principle that the total sum of the quantities being studied must be conserved by the system and its environment. Thus, the different phenomena that lead to transport are each considered individually with the knowledge that the sum of their contributions must equal zero. This principle is useful for calculating many relevant quantities. For example, in fluid mechanics, a common use of transport analysis is to determine the velocity profile of a fluid flowing through a rigid volume.
Transport phenomena are ubiquitous throughout the engineering disciplines. Some of the most common examples of transport analysis in engineering are seen in the fields of process, chemical, biological, and mechanical engineering, but the subject is a fundamental component of the curriculum in all disciplines involved in any way with fluid mechanics, heat transfer, and mass transfer. It is now considered to be a part of the engineering discipline as much as thermodynamics, mechanics, and electromagnetism.
Transport phenomena encompass all agents of physical change in the universe. Moreover, they are considered to be fundamental building blocks which developed the universe, and which are responsible for the success of all life on Earth. However, the scope here is limited to the relationship of transport phenomena to artificial engineered systems.
Overview
In physics, transport phenomena are all irreversible processes of statistical nature stemming from the random continuous motion of molecules, mostly observed in fluids. Every aspect of transport phenomena is grounded in two primary concepts : the conservation laws, and the constitutive equations. The conservation laws, which in the context of transport phenomena are formulated as continuity equations, describe how the quantity being studied must be conserved. The constitutive equations describe how the quantity in question responds to various stimuli via transport. Prominent examples include Fourier's law of heat conduction and the Navier–Stokes equations, which describe, respectively, the response of heat flux to temperature gradients and the relationship between fluid flux and the forces applied to the fluid. These equations also demonstrate the deep connection between transport phenomena and thermodynamics, a connection that explains why transport phenomena are irreversible. Almost all of these physical phenomena ultimately involve systems seeking their lowest energy state in keeping with the principle of minimum energy. As they approach this state, they tend to achieve true thermodynamic equilibrium, at which point there are no longer any driving forces in the system and transport ceases. The various aspects of such equilibrium are directly connected to a specific transport: heat transfer is the system's attempt to achieve thermal equilibrium with its environment, just as mass and momentum transport move the system towards chemical and mechanical equilibrium.
Examples of transport processes include heat conduction (energy transfer), fluid flow (momentum transfer), molecular diffusion (mass transfer), radiation and electric charge transfer in semiconductors.
Transport phenomena have wide application. For example, in solid state physics, the motion and interaction of electrons, holes and phonons are studied under "transport phenomena". Another example is in biomedical engineering, where some transport phenomena of interest are thermoregulation, perfusion, and microfluidics. In chemical engineering, transport phenomena are studied in reactor design, analysis of molecular or diffusive transport mechanisms, and metallurgy.
The transport of mass, energy, and momentum can be affected by the presence of external sources:
An odor dissipates more slowly (and may intensify) when the source of the odor remains present.
The rate of cooling of a solid that is conducting heat depends on whether a heat source is applied.
The gravitational force acting on a rain drop counteracts the resistance or drag imparted by the surrounding air.
Commonalities among phenomena
An important principle in the study of transport phenomena is analogy between phenomena.
Diffusion
There are some notable similarities in equations for momentum, energy, and mass transfer which can all be transported by diffusion, as illustrated by the following examples:
Mass: the spreading and dissipation of odors in air is an example of mass diffusion.
Energy: the conduction of heat in a solid material is an example of heat diffusion.
Momentum: the drag experienced by a rain drop as it falls in the atmosphere is an example of momentum diffusion (the rain drop loses momentum to the surrounding air through viscous stresses and decelerates).
The molecular transfer equations of Newton's law for fluid momentum, Fourier's law for heat, and Fick's law for mass are very similar. One can convert from one transport coefficient to another in order to compare all three different transport phenomena.
A great deal of effort has been devoted in the literature to developing analogies among these three transport processes for turbulent transfer so as to allow prediction of one from any of the others. The Reynolds analogy assumes that the turbulent diffusivities are all equal and that the molecular diffusivities of momentum (μ/ρ) and mass (DAB) are negligible compared to the turbulent diffusivities. When liquids are present and/or drag is present, the analogy is not valid. Other analogies, such as von Karman's and Prandtl's, usually result in poor relations.
The most successful and most widely used analogy is the Chilton and Colburn J-factor analogy. This analogy is based on experimental data for gases and liquids in both the laminar and turbulent regimes. Although it is based on experimental data, it can be shown to satisfy the exact solution derived from laminar flow over a flat plate. All of this information is used to predict transfer of mass.
Onsager reciprocal relations
In fluid systems described in terms of temperature, matter density, and pressure, it is known that temperature differences lead to heat flows from the warmer to the colder parts of the system; similarly, pressure differences will lead to matter flow from high-pressure to low-pressure regions (a "reciprocal relation"). What is remarkable is the observation that, when both pressure and temperature vary, temperature differences at constant pressure can cause matter flow (as in convection) and pressure differences at constant temperature can cause heat flow. The heat flow per unit of pressure difference and the density (matter) flow per unit of temperature difference are equal.
This equality was shown to be necessary by Lars Onsager using statistical mechanics as a consequence of the time reversibility of microscopic dynamics. The theory developed by Onsager is much more general than this example and capable of treating more than two thermodynamic forces at once.
Momentum transfer
In momentum transfer, the fluid is treated as a continuous distribution of matter. The study of momentum transfer, or fluid mechanics can be divided into two branches: fluid statics (fluids at rest), and fluid dynamics (fluids in motion).
When a fluid is flowing in the x-direction parallel to a solid surface, the fluid has x-directed momentum, and its concentration is υxρ. By random diffusion of molecules there is an exchange of molecules in the z-direction. Hence the x-directed momentum has been transferred in the z-direction from the faster- to the slower-moving layer.
The equation for momentum transfer is Newton's law of viscosity written as follows:
where τzx is the flux of x-directed momentum in the z-direction, ν is μ/ρ, the momentum diffusivity, z is the distance of transport or diffusion, ρ is the density, and μ is the dynamic viscosity. Newton's law of viscosity is the simplest relationship between the flux of momentum and the velocity gradient. It may be useful to note that this is an unconventional use of the symbol τzx; the indices are reversed as compared with standard usage in solid mechanics, and the sign is reversed.
Mass transfer
When a system contains two or more components whose concentration vary from point to point, there is a natural tendency for mass to be transferred, minimizing any concentration difference within the system. Mass transfer in a system is governed by Fick's first law: 'Diffusion flux from higher concentration to lower concentration is proportional to the gradient of the concentration of the substance and the diffusivity of the substance in the medium.' Mass transfer can take place due to different driving forces. Some of them are:
Mass can be transferred by the action of a pressure gradient (pressure diffusion)
Forced diffusion occurs because of the action of some external force
Diffusion can be caused by temperature gradients (thermal diffusion)
Diffusion can be caused by differences in chemical potential
This can be compared to Fick's law of diffusion, for a species A in a binary mixture consisting of A and B:
where D is the diffusivity constant.
Heat transfer
Many important engineered systems involve heat transfer. Some examples are the heating and cooling of process streams, phase changes, distillation, etc. The basic principle is the Fourier's law which is expressed as follows for a static system:
The net flux of heat through a system equals the conductivity times the rate of change of temperature with respect to position.
For convective transport involving turbulent flow, complex geometries, or difficult boundary conditions, the heat transfer may be represented by a heat transfer coefficient.
where A is the surface area, is the temperature driving force, Q is the heat flow per unit time, and h is the heat transfer coefficient.
Within heat transfer, two principal types of convection can occur:
Forced convection can occur in both laminar and turbulent flow. In the situation of laminar flow in circular tubes, several dimensionless numbers are used such as Nusselt number, Reynolds number, and Prandtl number. The commonly used equation is .
Natural or free convection is a function of Grashof and Prandtl numbers. The complexities of free convection heat transfer make it necessary to mainly use empirical relations from experimental data.
Heat transfer is analyzed in packed beds, nuclear reactors and heat exchangers.
Heat and mass transfer analogy
The heat and mass analogy allows solutions for mass transfer problems to be obtained from known solutions to heat transfer problems. Its arises from similar non-dimensional governing equations between heat and mass transfer.
Derivation
The non-dimensional energy equation for fluid flow in a boundary layer can simplify to the following, when heating from viscous dissipation and heat generation can be neglected:
Where and are the velocities in the x and y directions respectively normalized by the free stream velocity, and are the x and y coordinates non-dimensionalized by a relevant length scale, is the Reynolds number, is the Prandtl number, and is the non-dimensional temperature, which is defined by the local, minimum, and maximum temperatures:
The non-dimensional species transport equation for fluid flow in a boundary layer can be given as the following, assuming no bulk species generation:
Where is the non-dimensional concentration, and is the Schmidt number.
Transport of heat is driven by temperature differences, while transport of species is due to concentration differences. They differ by the relative diffusion of their transport compared to the diffusion of momentum. For heat, the comparison is between viscous diffusivity () and thermal diffusion (), given by the Prandtl number. Meanwhile, for mass transfer, the comparison is between viscous diffusivity () and mass Diffusivity (), given by the Schmidt number.
In some cases direct analytic solutions can be found from these equations for the Nusselt and Sherwood numbers. In cases where experimental results are used, one can assume these equations underlie the observed transport.
At an interface, the boundary conditions for both equations are also similar. For heat transfer at an interface, the no-slip condition allows us to equate conduction with convection, thus equating Fourier's law and Newton's law of cooling:
Where q” is the heat flux, is the thermal conductivity, is the heat transfer coefficient, and the subscripts and compare the surface and bulk values respectively.
For mass transfer at an interface, we can equate Fick's law with Newton's law for convection, yielding:
Where is the mass flux [kg/s ], is the diffusivity of species a in fluid b, and is the mass transfer coefficient. As we can see, and are analogous, and are analogous, while and are analogous.
Implementing the Analogy
Heat-Mass Analogy:
Because the Nu and Sh equations are derived from these analogous governing equations, one can directly swap the Nu and Sh and the Pr and Sc numbers to convert these equations between mass and heat.
In many situations, such as flow over a flat plate, the Nu and Sh numbers are functions of the Pr and Sc numbers to some coefficient . Therefore, one can directly calculate these numbers from one another using:
Where can be used in most cases, which comes from the analytical solution for the Nusselt Number for laminar flow over a flat plate. For best accuracy, n should be adjusted where correlations have a different exponent.
We can take this further by substituting into this equation the definitions of the heat transfer coefficient, mass transfer coefficient, and Lewis number, yielding:
For fully developed turbulent flow, with n=1/3, this becomes the Chilton–Colburn J-factor analogy. Said analogy also relates viscous forces and heat transfer, like the Reynolds analogy.
Limitations
The analogy between heat transfer and mass transfer is strictly limited to binary diffusion in dilute (ideal) solutions for which the mass transfer rates are low enough that mass transfer has no effect on the velocity field. The concentration of the diffusing species must be low enough that the chemical potential gradient is accurately represented by the concentration gradient (thus, the analogy has limited application to concentrated liquid solutions). When the rate of mass transfer is high or the concentration of the diffusing species is not low, corrections to the low-rate heat transfer coefficient can sometimes help. Further, in multicomponent mixtures, the transport of one species is affected by the chemical potential gradients of other species.
The heat and mass analogy may also break down in cases where the governing equations differ substantially. For instance, situations with substantial contributions from generation terms in the flow, such as bulk heat generation or bulk chemical reactions, may cause solutions to diverge.
Applications of the Heat-Mass Analogy
The analogy is useful for both using heat and mass transport to predict one another, or for understanding systems which experience simultaneous heat and mass transfer. For example, predicting heat transfer coefficients around turbine blades is challenging and is often done through measuring evaporating of a volatile compound and using the analogy. Many systems also experience simultaneous mass and heat transfer, and particularly common examples occur in processes with phase change, as the enthalpy of phase change often substantially influences heat transfer. Such examples include: evaporation at a water surface, transport of vapor in the air gap above a membrane distillation desalination membrane, and HVAC dehumidification equipment that combine heat transfer and selective membranes.
Applications
Pollution
The study of transport processes is relevant for understanding the release and distribution of pollutants into the environment. In particular, accurate modeling can inform mitigation strategies. Examples include the control of surface water pollution from urban runoff, and policies intended to reduce the copper content of vehicle brake pads in the U.S.
| Physical sciences | Physics basics: General | Physics |
27156495 | https://en.wikipedia.org/wiki/Lagerpetidae | Lagerpetidae | Lagerpetidae (; originally Lagerpetonidae) is a family of basal avemetatarsalians. Though traditionally considered the earliest-diverging dinosauromorphs (reptiles closer to dinosaurs than to pterosaurs), fossils described in 2020 suggest that lagerpetids may instead be pterosauromorphs (closer to pterosaurs). Lagerpetid fossils are known from the Triassic of San Juan (Argentina), Arizona (USA), Rio Grande do Sul (Brazil), Madagascar, New Mexico (USA), and Texas (USA). They were typically small, although some lagerpetids, like Dromomeron gigas and a specimen from the Santa Rosa Formation attributed to Dromomeron sp., were able to get quite large (femoral length ). Lagerpetid fossils are rare; the most common finds are bones of the hindlimbs, which possessed a number of unique features.
Description
As with most early avemetatarsalians, the most characteristic adaptations of lagerpetids occurred in their hip, leg and ankle bones, likely as a result of these being the bones most commonly preserved. Hip material is only known in Lagerpeton and Ixalerpeton, which share three adaptations of the ilium (upper blade of the hip). The supraacetabular crest, a ridge of bone which lies above the acetabulum (hip socket), is thickest above the middle portion of the acetabulum, rather than the front of it. However, it also extends further forwards than in most dinosauromorphs, snaking along the length of the pubic peduncle (the area of the ilium which connects to the pubis). The ilium's facet for the pubis opens downwards, a trait also acquired by ornithischian dinosaurs. The hip in general was wide, had a closed acetabulum (i.e. one with a bony inner wall), and had two sacral vertebrae, lacking many specializations of later dinosauromorphs, like dinosaurs.
Like other early archosaurs (and archosaur relatives such as Euparkeria), the femur (thigh bone) was slender and S-shaped. The femoral head was thin when seen from above, and its apex projected about 45 degrees between medially (inwards) and anteriorly (forwards). Most archosaurs had three tubera (bumps) on their flattened femoral head, one at the middle of the anterolateral (forwards/outwards) surface, another at the middle of the posteromedial (backwards/inwards) surface, and a small third one which was near the apex of the femoral head. However, lagerpetids lack the anterolateral tuber, instead having an emargination in the head just below where the tuber would normally be expected. The femoral head itself was notably hook-shaped when seen from the side. The distal portion of the femur (i.e. the portion near the knee) had a pair of condyles (knobs) on either side of the rear surface, as well as a third knob-like structure known as a crista tibiofibularis, which was present just above the lateral condyle. The crista tibiofibularis was uniquely enlarged in lagerpetids, and undergoes further evolution in Ixalerpeton and particularly Dromomeron.
The tibia and fibula (shin bones) were long and thin, with the tibia longer than the femur and generally resembling the tibia of early theropod dinosaurs. The ankle was formed by two main bones: the astragalus (which contacts both the tibia and fibula), and the calcaneum (which only contacts the fibula). As with dinosauromorphs, the astragalus was twice as wide as the reduced calcaneum. In addition, the two bones were co-ossified (fused together), akin to the condition in pterosaurs and some early dinosaurs (coelophysoids, for example). A pair of small, pyramid-shaped structures rise up out of the astragalus, one in front of the facet for the tibia, and the other behind it. The one in front is similar to a structure found in dinosauriform ankles known as the anterior ascending process, and it may be homologous with it. However, the posterior ascending process (the one behind the tibial facet) is entirely unique to lagerpetids. The rear of the astragalus lacks a horizontal groove, similar to Tropidosuchus, theropods, and ornithischians, but unlike most other archosauriforms. Like pterosaurs and dinosaurs (but unlike Marasuchus and most other archosaurs), the facet on the calcaneum which receives the fibula is concave and there is no evidence of a pronounced rearward bump known as a calcaneal tuber.
Classification
The lagerpetids were typically considered relatives of the dinosaurs, as a branch of the group Dinosauromorpha. The family was originally named Lagerpetonidae by Arcucci in 1986, though it was later renamed Lagerpetidae in a phylogenetic study by S. J. Nesbitt and colleagues in 2009. A clade of lagerpetids was also recovered in the large phylogenetic analyses of early dinosaurs and other dinosauromorphs that were produced by Baron, Norman & Barrett (2017). More recently, Muller et al. (2018) carried out the most comprehensive study on lagerpetid phylogeny, which assembled all lagerpetid specimens, taxa and morphotypes known so far into three of the most recent data matrices on early dinosauromorph/archosaur evolution. Finally, Garcia et al. (2019) added an unnamed lagerpetid (a new morphotype) to the data matrices used in the study by Muller et al. (2018).
Cladogram simplified after Cabreira et al., 2016:
By contrast, Kammerer et al. (2020), Ezcurra et al. (2020) recovered Lagerpetidae as the sister clade to pterosaurs, based on newly-described fossils of the jaw, forelimbs, and braincase. Baron (2021) also recovered a similar result.
| Biology and health sciences | Other prehistoric archosaurs | Animals |
27158894 | https://en.wikipedia.org/wiki/Addiction | Addiction | Addiction is a neuropsychological disorder characterized by a persistent and intense urge to use a drug or engage in a behavior that produces natural reward, despite substantial harm and other negative consequences. Repetitive drug use can alter brain function in synapses similar to natural rewards like food or falling in love in ways that perpetuate craving and weakens self-control for people with pre-existing vulnerabilities. This phenomenon – drugs reshaping brain function – has led to an understanding of addiction as a brain disorder with a complex variety of psychosocial as well as neurobiological factors that are implicated in the development of addiction. While mice given cocaine showed the compulsive and involuntary nature of addiction, for humans this is more complex, related to behavior or personality traits.
Classic signs of addiction include compulsive engagement in rewarding stimuli, preoccupation with substances or behavior, and continued use despite negative consequences. Habits and patterns associated with addiction are typically characterized by immediate gratification (short-term reward), coupled with delayed deleterious effects (long-term costs).
Examples of substance addiction include alcoholism, cannabis addiction, amphetamine addiction, cocaine addiction, nicotine addiction, opioid addiction, and eating or food addiction. Behavioral addictions may include gambling addiction, shopping addiction, stalking, internet addiction, social media addiction, video game addiction, and sexual addiction. The DSM-5 and ICD-10 only recognize gambling addictions as behavioral addictions, but the ICD-11 also recognizes gaming addictions.
Signs and symptoms
Signs and symptoms of drug addiction can vary depending on the type of addiction. Symptoms may include:
Continuation of drug use despite the knowledge of consequences
Disregarding financial status when it comes to drug purchases
Ensuring a stable supply of the drug
Needing more of the drug over time to achieve similar effects
Social and work life impacted due to drug use
Unsuccessful attempts to stop drug use
Urge to use drug regularly
Other signs and symptoms can be categorized across relevant dimensions:
Definitions
The word "addiction" derives from the Latin "addico", meaning "giving over" with both positive connotations (devotion, dedication) and negative ones (being enslaved to a creditor in Roman law). This dual meaning persisted in traditional English dictionaries, encompassing both legal surrender and personal devotion to habits. Later, 19th century temperance movements narrowed the definition of addiction to just drug-related disease, ignoring behavioral addictions and the possibility of positive or neutral addictions. This restrictive view opposes the current understanding of addiction.
"Addiction" and "addictive behaviour" are polysemes denoting a category of mental disorders, of neuropsychological symptoms, or of merely maladaptive/harmful habits and lifestyles. A common use of "addiction" in medicine is for neuropsychological symptoms denoting pervasive/excessive and intense urges to engage in a category of behavioral compulsions or impulses towards sensory rewards (e.g., alcohol, betel quid, drugs, sex, gambling, video gaming). Addictive disorders or addiction disorders are mental disorders involving high intensities of addictions (as neuropsychological symptoms) that induce functional disabilities (i.e., limit subjects' social/family and occupational activities); the two categories of such disorders are substance-use addictions and behavioral addictions.
The DSM-5 classifies addiction the most severe stage of substance use disorder, due to significant loss of control and the presence of compulsive behaviours despite the desire to stop. It is a definition that many scientific papers and reports use.
"Dependence" is also a polyseme denoting either neuropsychological symptoms or mental disorders. In the DSM-5, dependences differ from addictions and can even normally happen without addictions; besides, substance-use dependences are severe stages of substance-use addictions (i.e. mental disorders) involving withdrawal issues. In the ICD-11, "substance-use dependence" is a synonym of "substance-use addiction" (i.e. neuropsychological symptoms) that can but do not necessarily involve withdrawal issues.
Substance addiction
Drug addiction
Drug addiction, which belongs to the class of substance-related disorders, is a chronic and relapsing brain disorder that features drug seeking and drug abuse, despite their harmful effects. This form of addiction changes brain circuitry such that the brain's reward system is compromised, causing functional consequences for stress management and self-control. Damage to the functions of the organs involved can persist throughout a lifetime and cause death if untreated. Substances involved with drug addiction include alcohol, nicotine, marijuana, opioids, cocaine, amphetamines, and even foods with high fat and sugar content. Addictions can begin experimentally in social contexts and can arise from the use of prescribed medications or a variety of other measures.
Drug addiction has been shown to work in phenomenological, conditioning (operant and classical), cognitive models, and the cue reactivity model. However, no one model completely illustrates substance abuse.
Risk factors for addiction include:
Aggressive behavior (particularly in childhood)
Availability of substance
Community economic status
Experimentation
Epigenetics
Impulsivity (attentional, motor, or non-planning)
Lack of parental supervision
Lack of peer refusal skills
Mental disorders
Method substance is taken
Usage of substance in youth
Food addiction
The diagnostic criteria for food or eating addiction has not been categorized or defined in references such as the Diagnostic and Statistical Manual of Mental Disorders (DSM or DSM-5) and is based on subjective experiences similar to substance use disorders. Food addiction may be found in those with eating disorders, though not all people with eating disorders have food addiction and not all of those with food addiction have a diagnosed eating disorder. Long-term frequent and excessive consumption of foods high in fat, salt, or sugar, such as chocolate, can produce an addiction similar to drugs since they trigger the brain's reward system, such that the individual may desire the same foods to an increasing degree over time. The signals sent when consuming highly palatable foods have the ability to counteract the body's signals for fullness and persistent cravings will result. Those who show signs of food addiction may develop food tolerances, in which they eat more, despite the food becoming less satisfactory.
Chocolate's sweet flavor and pharmacological ingredients are known to create a strong craving or feel 'addictive' by the consumer. A person who has a strong liking for chocolate may refer to themselves as a chocoholic.
Risk factors for developing food addiction include excessive overeating and impulsivity.
The Yale Food Addiction Scale (YFAS), version 2.0, is the current standard measure for assessing whether an individual exhibits signs and symptoms of food addiction. It was developed in 2009 at Yale University on the hypothesis that foods high in fat, sugar, and salt have addictive-like effects which contribute to problematic eating habits. The YFAS is designed to address 11 substance-related and addictive disorders (SRADs) using a 25-item self-report questionnaire, based on the diagnostic criteria for SRADs as per DSM-5. A potential food addiction diagnosis is predicted by the presence of at least two out of 11 SRADs and a significant impairment to daily activities.
The Barratt Impulsiveness Scale, specifically the BIS-11 scale, and the UPPS-P Impulsive Behavior subscales of Negative Urgency and Lack of Perseverance have been shown to have relation to food addiction.
Behavioral addiction
The term behavioral addiction refers to a compulsion to engage in a natural reward – which is a behavior that is inherently rewarding (i.e., desirable or appealing) – despite adverse consequences. Preclinical evidence has demonstrated that marked increases in the expression of ΔFosB through repetitive and excessive exposure to a natural reward induces the same behavioral effects and neuroplasticity as occurs in a drug addiction.
Addiction can exist without psychotropic drugs, an idea that was popularized by psychologist Stanton Peele. These are termed behavioral addictions. Such addictions may be passive or active, but they commonly contain reinforcing features, which are found in most addictions. Sexual behavior, eating, gambling, playing video games, and shopping are all associated with compulsive behaviors in humans and have been shown to activate the mesolimbic pathway and other parts of the reward system. Based on this evidence, sexual addiction, gambling addiction, video game addiction, and shopping addiction are classified accordingly.
Sexual
Gambling
Gambling provides a natural reward that is associated with compulsive behavior. Functional neuroimaging evidence shows that gambling activates the reward system and the mesolimbic pathway in particular. It is known that dopamine is involved in learning, motivation, as well as the reward system. The exact role of dopamine in gambling addiction has been debated. Suggested roles for D2, D3, and D4 dopamine receptors, as well as D3 receptors in the substantia nigra have been found in rat and human models, showing a correlation with the severity of the gambling behavior. This in turn was linked with greater dopamine release in the dorsal striatum.
Gambling addictions are linked with comorbidities such as mental health disorders, substance abuse, alcohol use disorder, and personality disorders.
Risk factors for gambling addictions include antisocial behavior, impulsive personality, male sex, sensation seeking, substance use, and young age.
Gambling addiction has been associated with some personality traits, including: harm avoidance, low self direction, decision making and planning insufficiencies, impulsivity, as well as sensation seeking individuals. Although some personality traits can be linked with gambling addiction, there is no general description of individuals addicted to gambling.
Internet
Internet addiction does not have any standardized definition, yet there is widespread agreement that this problem exists. Debate over the classification of problematic internet use considers whether it should be thought of as a behavioral addiction, an impulse control disorder, or an obsessive-compulsive disorder. Others argue that internet addiction should be considered a symptom of an underlying mental health condition and not a disorder in itself. Internet addiction has been described as "a psychological dependence on the Internet, regardless of the type of activity once logged on." Problematic internet use may include a preoccupation with the internet and/or digital media, excessive time spent using the internet despite resultant distress in the individual, increase in the amount of internet use required to achieve the same desired emotional response, loss of control over one's internet use habits, withdrawal symptoms, and continued problematic internet use despite negative consequences to one's work, social, academic, or personal life.
Studies conducted in India, United States, Asia, and Europe have identified Internet addiction prevalence rates ranging in value from 1% to 19%, with the adolescent population having high rates compared to other age groups. Prevalence rates have been difficult to establish due to a lack of universally accepted diagnostic criteria, a lack of diagnostic instruments demonstrating cross-cultural validity and reliability, and existing controversy surrounding the validity of labeling problematic internet use as an addictive disorder. The most common scale used to measure addiction is the Internet Addiction Test developed by Kimberly Young.
People with internet addiction are likely to have a comorbid psychiatric disorder. Comorbid diagnoses identified alongside internet addiction include affective mood disorders, anxiety disorders, substance use disorders, and attention deficit hyperactivity disorder.
Video games
Video game addiction is characterized by the World Health Organization (WHO) as excessive gaming behavior, potentially prioritized over other interests, despite the negative consequences that may arise, for a period of at least 12 months. In May 2019, the WHO introduced gaming disorder in the 11th edition of the International Classification of Diseases. Video game addiction has been shown to be more prevalent in males than females, higher by 2.9 times. It has been suggested that people of younger ages are more prone to become addicted to video games. People with certain personalities may be more susceptible to gaming addictions.
Risk factors for video game addiction include:
Male,
Psychopathologies (e.g. ADHD or MDD), and
Social anxiety.
Shopping
Shopping addiction, or compulsive buying disorder (CBD), is the excessive urge to shop or spend, potentially resulting in unwanted consequences. These consequences can have serious impacts, such as increased consumer debt, negatively affected relationships, increased risk of illegal behavior, and suicide attempts. Shopping addiction occurs worldwide and has shown a 5.8% prevalence in the United States. Similar to other behavioral addictions, CBD can be linked to mood disorders, substance use disorders, eating disorders, and other disorders involving a lack of control.
Screening and assessment
Addictions Neuroclinical Assessment
The Addictions Neuroclinical Assessment is used to diagnose addiction disorders. This tool measures three different domains: executive function, incentive salience, and negative emotionality. Executive functioning consists of processes that would be disrupted in addiction. In the context of addiction, incentive salience determines how one perceives the addictive substance. Increased negative emotional responses have been found with individuals with addictions.
Tobacco, Alcohol, Prescription Medication, and Other Substance Use (TAPS)
This is a screening and assessment tool in one, assessing commonly used substances. This tool allows for a simple diagnosis, eliminating the need for several screening and assessment tools, as it includes both TAPS-1 and TAPS-2, screening and assessment tools respectively. The screening component asks about the frequency of use of the specific substance (tobacco, alcohol, prescription medication, and other). If an individual screens positive, the second component will begin. This dictates the risk level of the substance.
CRAFFT
The CRAFFT (Car-Relax-Alone-Forget-Family and Friends-Trouble) is a screening tool that is used in medical centers. The CRAFFT is in version 2.1 and has a version for nicotine and tobacco use called the CRAFFT 2.1+N. This tool is used to identify substance use, substance related driving risk, and addictions among adolescents. This tool uses a set of questions for different scenarios. In the case of a specific combination of answers, different question sets can be used to yield a more accurate answer. After the questions, the DSM-5 criteria are used to identify the likelihood of the person having substance use disorder. After these tests are done, the clinician is to give the "5 RS" of brief counseling.
The five Rs of brief counseling includes:
REVIEW screening results
RECOMMEND to not use
RIDING/DRIVING risk counseling
RESPONSE: elicit self-motivational statements
REINFORCE self-efficacy
Drug Abuse Screening Test (DAST-10)
The Drug Abuse Screening Test (DAST) is a self-reporting tool that measures problematic substance use. Responses to this test are recorded as yes or no answers, and scored as a number between zero and 28. Drug abuse or dependence, are indicated by a cut off score of 6. Three versions of this screening tool are in use: DAST-28, DAST-20, and DAST-10. Each of these instruments are copyrighted by Dr. Harvey A. Skinner.
Alcohol, Smoking, and Substance Involvement Test (ASSIST)
The Alcohol, Smoking, and Substance Involvement Test (ASSIST) is an interview-based questionnaire consisting of eight questions developed by the WHO. The questions ask about lifetime use; frequency of use; urge to use; frequency of health, financial, social, or legal problems related to use; failure to perform duties; if anyone has raised concerns over use; attempts to limit or moderate use; and use by injection.
Causes
Personality theories
Personality theories of addiction are psychological models that associate personality traits or modes of thinking (i.e., affective states) with an individual's proclivity for developing an addiction. Data analysis demonstrates that psychological profiles of drug users and non-users have significant differences and the psychological predisposition to using different drugs may be different. Models of addiction risk that have been proposed in psychology literature include: an affect dysregulation model of positive and negative psychological affects, the reinforcement sensitivity theory of impulsiveness and behavioral inhibition, and an impulsivity model of reward sensitization and impulsiveness.
Neuropsychology
The transtheoretical model of change (TTM) can point to how someone may be conceptualizing their addiction and the thoughts around it, including not being aware of their addiction.
Cognitive control and stimulus control, which is associated with operant and classical conditioning, represent opposite processes (i.e., internal vs external or environmental, respectively) that compete over the control of an individual's elicited behaviors. Cognitive control, and particularly inhibitory control over behavior, is impaired in both addiction and attention deficit hyperactivity disorder. Stimulus-driven behavioral responses (i.e., stimulus control) that are associated with a particular rewarding stimulus tend to dominate one's behavior in an addiction.
Stimulus control of behavior
In operant conditioning, behavior is influenced by outside stimulus, such as a drug. The operant conditioning theory of learning is useful in understanding why the mood-altering or stimulating consequences of drug use can reinforce continued use (an example of positive reinforcement) and why the addicted person seeks to avoid withdrawal through continued use (an example of negative reinforcement). Stimulus control is using the absence of the stimulus or presence of a reward to influence the resulting behavior.
Cognitive control of behavior
Cognitive control is the intentional selection of thoughts, behaviors, and emotions, based on our environment. It has been shown that drugs alter the way our brains function, and its structure. Cognitive functions such as learning, memory, and impulse control, are affected by drugs. These effects promote drug use, as well as hinder the ability to abstain from it. The increase in dopamine release is prominent in drug use, specifically in the ventral striatum and the nucleus accumbens. Dopamine is responsible for producing pleasurable feelings, as well driving us to perform important life activities. Addictive drugs cause a significant increase in this reward system, causing a large increase in dopamine signaling as well as increase in reward-seeking behavior, in turn motivating drug use. This promotes the development of a maladaptive drug to stimulus relationship. Early drug use leads to these maladaptive associations, later affecting cognitive processes used for coping, which are needed to successfully abstain from them.
Risk factors
A number of genetic and environmental risk factors exist for developing an addiction. Genetic and environmental risk factors each account for roughly half of an individual's risk for developing an addiction; the contribution from epigenetic risk factors to the total risk is unknown. Even in individuals with a relatively low genetic risk, exposure to sufficiently high doses of an addictive drug for a long period of time (e.g., weeks–months) can result in an addiction. Adverse childhood events are associated with negative health outcomes, such as substance use disorder. Childhood abuse or exposure to violent crime is related to developing a mood or anxiety disorder, as well as a substance dependence risk.
Genetic factors
Genetic factors, along with socio-environmental (e.g., psychosocial) factors, have been established as significant contributors to addiction vulnerability. Studies done on 350 hospitalized drug-dependent patients showed that over half met the criteria for alcohol abuse, with a role of familial factors being prevalent. Genetic factors account for 40–60% of the risk factors for alcoholism. Similar rates of heritability for other types of drug addiction have been indicated, specifically in genes that encode the Alpha5 Nicotinic Acetylcholine Receptor. Knestler hypothesized in 1964 that a gene or group of genes might contribute to predisposition to addiction in several ways. For example, altered levels of a normal protein due to environmental factors may change the structure or functioning of specific brain neurons during development. These altered brain neurons could affect the susceptibility of an individual to an initial drug use experience. In support of this hypothesis, animal studies have shown that environmental factors such as stress can affect an animal's genetic expression.
In humans, twin studies into addiction have provided some of the highest-quality evidence of this link, with results finding that if one twin is affected by addiction, the other twin is likely to be as well, and to the same substance. Further evidence of a genetic component is research findings from family studies which suggest that if one family member has a history of addiction, the chances of a relative or close family developing those same habits are much higher than one who has not been introduced to addiction at a young age.
The data implicating specific genes in the development of drug addiction is mixed for most genes. Many addiction studies that aim to identify specific genes focus on common variants with an allele frequency of greater than 5% in the general population. When associated with disease, these only confer a small amount of additional risk with an odds ratio of 1.1–1.3 percent; this has led to the development the rare variant hypothesis, which states that genes with low frequencies in the population (<1%) confer much greater additional risk in the development of the disease.
Genome-wide association studies (GWAS) are used to examine genetic associations with dependence, addiction, and drug use. These studies rarely identify genes from proteins previously described via animal knockout models and candidate gene analysis. Instead, large percentages of genes involved in processes such as cell adhesion are commonly identified. The important effects of endophenotypes are typically not capable of being captured by these methods. Genes identified in GWAS for drug addiction may be involved either in adjusting brain behavior before drug experiences, subsequent to them, or both.
Environmental factors
Environmental risk factors for addiction are the experiences of an individual during their lifetime that interact with the individual's genetic composition to increase or decrease his or her vulnerability to addiction. For example, after the nationwide outbreak of COVID-19, more people quit (vs. started) smoking; and smokers, on average, reduced the quantity of cigarettes they consumed. More generally, a number of different environmental factors have been implicated as risk factors for addiction, including various psychosocial stressors. The National Institute on Drug Abuse (NIDA) and studies cite lack of parental supervision, the prevalence of peer substance use, substance availability, and poverty as risk factors for substance use among children and adolescents. The brain disease model of addiction posits that an individual's exposure to an addictive drug is the most significant environmental risk factor for addiction. Many researchers, including neuroscientists, indicate that the brain disease model presents a misleading, incomplete, and potentially detrimental explanation of addiction.
The psychoanalytic theory model defines addiction as a form of defense against feelings of hopelessness and helplessness as well as a symptom of failure to regulate powerful emotions related to adverse childhood experiences (ACEs), various forms of maltreatment and dysfunction experienced in childhood. In this case, the addictive substance provides brief but total relief and positive feelings of control. The Adverse Childhood Experiences Study by the Centers for Disease Control and Prevention has shown a strong dose–response relationship between ACEs and numerous health, social, and behavioral problems throughout a person's lifespan, including substance use disorder. Children's neurological development can be permanently disrupted when they are chronically exposed to stressful events such as physical, emotional, or sexual abuse, physical or emotional neglect, witnessing violence in the household, or a parent being incarcerated or having a mental illness. As a result, the child's cognitive functioning or ability to cope with negative or disruptive emotions may be impaired. Over time, the child may adopt substance use as a coping mechanism or as a result of reduced impulse control, particularly during adolescence. Vast amounts of children who experienced abuse have gone on to have some form of addiction in their adolescence or adult life. This pathway towards addiction that is opened through stressful experiences during childhood can be avoided by a change in environmental factors throughout an individual's life and opportunities of professional help. If one has friends or peers who engage in drug use favorably, the chances of them developing an addiction increases. Family conflict and home management is a cause for one to become engaged in drug use.
Social control theory
According to Travis Hirschi's social control theory, adolescents with stronger attachments to family, religious, academic, and other social institutions are less likely to engage in delinquent and maladaptive behavior such as drug use leading to addiction.
Age
Adolescence represents a period of increased vulnerability for developing an addiction. In adolescence, the incentive-rewards systems in the brain mature well before the cognitive control center. This consequentially grants the incentive-rewards systems a disproportionate amount of power in the behavioral decision-making process. Therefore, adolescents are increasingly likely to act on their impulses and engage in risky, potentially addicting behavior before considering the consequences. Not only are adolescents more likely to initiate and maintain drug use, but once addicted they are more resistant to treatment and more liable to relapse.
Most individuals are exposed to and use addictive drugs for the first time during their teenage years. In the United States, there were just over 2.8 million new users of illicit drugs in 2013 (7,800 new users per day); among them, 54.1% were under 18 years of age. In 2011, there were approximately 20.6 million people in the United States over the age of 12 with an addiction. Over 90% of those with an addiction began drinking, smoking or using illicit drugs before the age of 18.
Comorbid disorders
Individuals with comorbid (i.e., co-occurring) mental health disorders such as depression, anxiety, attention-deficit/hyperactivity disorder (ADHD) or post-traumatic stress disorder are more likely to develop substance use disorders. The cites early aggressive behavior as a risk factor for substance use. The National Bureau of Economic Research found that there is a "definite connection between mental illness and the use of addictive substances" and a majority of mental health patients participate in the use of these substances: 38% alcohol, 44% cocaine, and 40% cigarettes.
Epigenetic
Epigenetics is the study of stable phenotypic changes that do not involve alterations in the DNA sequence. Illicit drug use has been found to cause epigenetic changes in DNA methylation, as well as chromatin remodeling. The epigenetic state of chromatin may pose as a risk for the development of substance addictions. It has been found that emotional stressors, as well as social adversities may lead to an initial epigenetic response, which causes an alteration to the reward-signalling pathways. This change may predispose one to experience a positive response to drug use.
Transgenerational epigenetic inheritance
Epigenetic genes and their products (e.g., proteins) are the key components through which environmental influences can affect the genes of an individual: they serve as the mechanism responsible for transgenerational epigenetic inheritance, a phenomenon in which environmental influences on the genes of a parent can affect the associated traits and behavioral phenotypes of their offspring (e.g., behavioral responses to environmental stimuli). In addiction, epigenetic mechanisms play a central role in the pathophysiology of the disease; it has been noted that some of the alterations to the epigenome which arise through chronic exposure to addictive stimuli during an addiction can be transmitted across generations, in turn affecting the behavior of one's children (e.g., the child's behavioral responses to addictive drugs and natural rewards).
The general classes of epigenetic alterations that have been implicated in transgenerational epigenetic inheritance include DNA methylation, histone modifications, and downregulation or upregulation of microRNAs. With respect to addiction, more research is needed to determine the specific heritable epigenetic alterations that arise from various forms of addiction in humans and the corresponding behavioral phenotypes from these epigenetic alterations that occur in human offspring. Based on preclinical evidence from animal research, certain addiction-induced epigenetic alterations in rats can be transmitted from parent to offspring and produce behavioral phenotypes that decrease the offspring's risk of developing an addiction. More generally, the heritable behavioral phenotypes that are derived from addiction-induced epigenetic alterations and transmitted from parent to offspring may serve to either increase or decrease the offspring's risk of developing an addiction.
Mechanisms
Addiction is a disorder of the brain's reward system developing through transcriptional and epigenetic mechanisms as a result of chronically high levels of exposure to an addictive stimulus (e.g., eating food, the use of cocaine, engagement in sexual activity, participation in high-thrill cultural activities such as gambling, etc.) over extended time. DeltaFosB (ΔFosB), a gene transcription factor, is a critical component and common factor in the development of virtually all forms of behavioral and drug addictions. Two decades of research into ΔFosB's role in addiction have demonstrated that addiction arises, and the associated compulsive behavior intensifies or attenuates, along with the overexpression of ΔFosB in the D1-type medium spiny neurons of the nucleus accumbens. Due to the causal relationship between ΔFosB expression and addictions, it is used preclinically as an addiction biomarker. ΔFosB expression in these neurons directly and positively regulates drug self-administration and reward sensitization through positive reinforcement, while decreasing sensitivity to aversion.
Chronic addictive drug use causes alterations in gene expression in the mesocorticolimbic projection. The most important transcription factors that produce these alterations are ΔFosB, cAMP response element binding protein (CREB), and nuclear factor kappa B (NF-κB). ΔFosB is the most significant biomolecular mechanism in addiction because the overexpression of ΔFosB in the D1-type medium spiny neurons in the nucleus accumbens is necessary and sufficient for many of the neural adaptations and behavioral effects (e.g., expression-dependent increases in drug self-administration and reward sensitization) seen in drug addiction. ΔFosB expression in nucleus accumbens D1-type medium spiny neurons directly and positively regulates drug self-administration and reward sensitization through positive reinforcement while decreasing sensitivity to aversion. ΔFosB has been implicated in mediating addictions to many different drugs and drug classes, including alcohol, amphetamine and other substituted amphetamines, cannabinoids, cocaine, methylphenidate, nicotine, opiates, phenylcyclidine, and propofol, among others. ΔJunD, a transcription factor, and G9a, a histone methyltransferase, both oppose the function of ΔFosB and inhibit increases in its expression. Increases in nucleus accumbens ΔJunD expression (via viral vector-mediated gene transfer) or G9a expression (via pharmacological means) reduces, or with a large increase can even block, many of the neural and behavioral alterations that result from chronic high-dose use of addictive drugs (i.e., the alterations mediated by ΔFosB).
ΔFosB plays an important role in regulating behavioral responses to natural rewards, such as palatable food, sex, and exercise. Natural rewards, like drugs of abuse, induce gene expression of ΔFosB in the nucleus accumbens, and chronic acquisition of these rewards can result in a similar pathological addictive state through ΔFosB overexpression. Consequently, ΔFosB is the key transcription factor involved in addictions to natural rewards (i.e., behavioral addictions) as well; in particular, ΔFosB in the nucleus accumbens is critical for the reinforcing effects of sexual reward. Research on the interaction between natural and drug rewards suggests that dopaminergic psychostimulants (e.g., amphetamine) and sexual behavior act on similar biomolecular mechanisms to induce ΔFosB in the nucleus accumbens and possess bidirectional cross-sensitization effects that are mediated through ΔFosB. This phenomenon is notable since, in humans, a dopamine dysregulation syndrome, characterized by drug-induced compulsive engagement in natural rewards (specifically, sexual activity, shopping, and gambling), has been observed in some individuals taking dopaminergic medications.
ΔFosB inhibitors (drugs or treatments that oppose its action) may be an effective treatment for addiction and addictive disorders.
The release of dopamine in the nucleus accumbens plays a role in the reinforcing qualities of many forms of stimuli, including naturally reinforcing stimuli like palatable food and sex. Altered dopamine neurotransmission is frequently observed following the development of an addictive state. In humans and lab animals that have developed an addiction, alterations in dopamine or opioid neurotransmission in the nucleus accumbens and other parts of the striatum are evident. Use of certain drugs (e.g., cocaine) affect cholinergic neurons that innervate the reward system, in turn affecting dopamine signaling in this region.
A recent study in Addiction reports that GLP-1 agonist medications, such as semaglutide, which are commonly used for diabetes and weight management, may also reduce the risk of overdose and alcohol intoxication in people with substance use disorders. The study analyzed nearly nine years of health records from 1.3 million individuals across 136 U.S. hospitals, including 500,000 with opioid use disorder and over 800,000 with alcohol use disorder. Researchers found that those who used Ozempic or similar medications had a 40% lower risk of opioid overdose and a 50% lower risk of alcohol intoxication compared to those not using these drugs.
Reward system
Mesocorticolimbic pathway
Understanding the pathways in which drugs act and how drugs can alter those pathways is key when examining the biological basis of drug addiction. The reward pathway, known as the mesolimbic pathway, or its extension, the mesocorticolimbic pathway, is characterized by the interaction of several areas of the brain.
The projections from the ventral tegmental area (VTA) are a network of dopaminergic neurons with co-localized postsynaptic glutamate receptors (AMPAR and NMDAR). These cells respond when stimuli indicative of a reward are present. The VTA supports learning and sensitization development and releases dopamine (DA) into the forebrain. These neurons project and release DA into the nucleus accumbens, through the mesolimbic pathway. Virtually all drugs causing drug addiction increase the DA release in the mesolimbic pathway.
The nucleus accumbens (NAcc) is one output of the VTA projections. The nucleus accumbens itself consists mainly of GABAergic medium spiny neurons (MSNs). The NAcc is associated with acquiring and eliciting conditioned behaviors, and is involved in the increased sensitivity to drugs as addiction progresses. Overexpression of ΔFosB in the nucleus accumbens is a necessary common factor in essentially all known forms of addiction; ΔFosB is a strong positive modulator of positively reinforced behaviors.
The prefrontal cortex, including the anterior cingulate and orbitofrontal cortices, is another VTA output in the mesocorticolimbic pathway; it is important for the integration of information which helps determine whether a behavior will be elicited. It is critical for forming associations between the rewarding experience of drug use and cues in the environment. Importantly, these cues are strong mediators of drug-seeking behavior and can trigger relapse even after months or years of abstinence.
Other brain structures that are involved in addiction include:
The basolateral amygdala projects into the NAcc and is thought to be important for motivation.
The hippocampus is involved in drug addiction, because of its role in learning and memory. Much of this evidence stems from investigations showing that manipulating cells in the hippocampus alters DA levels in NAcc and firing rates of VTA dopaminergic cells.
Role of dopamine and glutamate
Dopamine is the primary neurotransmitter of the reward system in the brain. It plays a role in regulating movement, emotion, cognition, motivation, and feelings of pleasure. Natural rewards, like eating, as well as recreational drug use cause a release of dopamine, and are associated with the reinforcing nature of these stimuli. Nearly all addictive drugs, directly or indirectly, act on the brain's reward system by heightening dopaminergic activity.
Excessive intake of many types of addictive drugs results in repeated release of high amounts of dopamine, which in turn affects the reward pathway directly through heightened dopamine receptor activation. Prolonged and abnormally high levels of dopamine in the synaptic cleft can induce receptor downregulation in the neural pathway. Downregulation of mesolimbic dopamine receptors can result in a decrease in the sensitivity to natural reinforcers.
Drug seeking behavior is induced by glutamatergic projections from the prefrontal cortex to the nucleus accumbens. This idea is supported with data from experiments showing that drug seeking behavior can be prevented following the inhibition of AMPA glutamate receptors and glutamate release in the nucleus accumbens.
Reward sensitization
Reward sensitization is a process that causes an increase in the amount of reward (specifically, incentive salience) that is assigned by the brain to a rewarding stimulus (e.g., a drug). In simple terms, when reward sensitization to a specific stimulus (e.g., a drug) occurs, an individual's "wanting" or desire for the stimulus itself and its associated cues increases. Reward sensitization normally occurs following chronically high levels of exposure to the stimulus. ΔFosB expression in D1-type medium spiny neurons in the nucleus accumbens has been shown to directly and positively regulate reward sensitization involving drugs and natural rewards.
"Cue-induced wanting" or "cue-triggered wanting", a form of craving that occurs in addiction, is responsible for most of the compulsive behavior that people with addictions exhibit. During the development of an addiction, the repeated association of otherwise neutral and even non-rewarding stimuli with drug consumption triggers an associative learning process that causes these previously neutral stimuli to act as conditioned positive reinforcers of addictive drug use (i.e., these stimuli start to function as drug cues). As conditioned positive reinforcers of drug use, these previously neutral stimuli are assigned incentive salience (which manifests as a craving) – sometimes at pathologically high levels due to reward sensitization – which can transfer to the primary reinforcer (e.g., the use of an addictive drug) with which it was originally paired.
Research on the interaction between natural and drug rewards suggests that dopaminergic psychostimulants (e.g., amphetamine) and sexual behavior act on similar biomolecular mechanisms to induce ΔFosB in the nucleus accumbens and possess a bidirectional reward cross-sensitization effect that is mediated through ΔFosB. In contrast to ΔFosB's reward-sensitizing effect, CREB transcriptional activity decreases user's sensitivity to the rewarding effects of the substance. CREB transcription in the nucleus accumbens is implicated in psychological dependence and symptoms involving a lack of pleasure or motivation during drug withdrawal.
Neuroepigenetic mechanisms
Altered epigenetic regulation of gene expression within the brain's reward system plays a significant and complex role in the development of drug addiction. Addictive drugs are associated with three types of epigenetic modifications within neurons. These are (1) histone modifications, (2) epigenetic methylation of DNA at CpG sites at (or adjacent to) particular genes, and (3) epigenetic downregulation or upregulation of microRNAs which have particular target genes. As an example, while hundreds of genes in the cells of the nucleus accumbens (NAc) exhibit histone modifications following drug exposure – particularly, altered acetylation and methylation states of histone residues – most other genes in the NAc cells do not show such changes.
Diagnosis
Classification
DSM-5
The fifth edition of the DSM uses the term substance use disorder to refer to a spectrum of drug use-related disorders. The DSM-5 eliminates the terms abuse and dependence from diagnostic categories, instead using the specifiers of mild, moderate and severe to indicate the extent of disordered use. These specifiers are determined by the number of diagnostic criteria present in a given case. In the DSM-5, the term drug addiction is synonymous with severe substance use disorder.
The DSM-5 introduced a new diagnostic category for behavioral addictions. Problem gambling is the only condition included in this category in the fifth edition. Internet gaming disorder is listed as a "condition requiring further study" in the DSM-5.
Past editions have used physical dependence and the associated withdrawal syndrome to identify an addictive state. Physical dependence occurs when the body has adjusted by incorporating the substance into its "normal" functioning – i.e., attains homeostasis – and therefore physical withdrawal symptoms occur on cessation of use. Tolerance is the process by which the body continually adapts to the substance and requires increasingly larger amounts to achieve the original effects. Withdrawal refers to physical and psychological symptoms experienced when reducing or discontinuing a substance that the body has become dependent on. Symptoms of withdrawal generally include but are not limited to body aches, anxiety, irritability, intense cravings for the substance, dysphoria, nausea, hallucinations, headaches, cold sweats, tremors, and seizures. During acute physical opioid withdrawal, symptoms of restless legs syndrome are common and may be profound. This phenomenon originated the idiom "kicking the habit".
Medical researchers who actively study addiction have criticized the DSM classification of addiction for being flawed and involving arbitrary diagnostic criteria.
ICD-11
The eleventh revision of the International Classification of Diseases, commonly referred to as ICD-11, conceptualizes diagnosis somewhat differently. ICD-11 first distinguishes between problems with psychoactive substance use ("Disorders due to substance use") and behavioral addictions ("Disorders due to addictive behaviours"). With regard to psychoactive substances, ICD-11 explains that the included substances initially produce "pleasant or appealing psychoactive effects that are rewarding and reinforcing with repeated use, [but] with continued use, many of the included substances have the capacity to produce dependence. They have the potential to cause numerous forms of harm, both to mental and physical health." Instead of the DSM-5 approach of one diagnosis ("Substance Use Disorder") covering all types of problematic substance use, ICD-11 offers three diagnostic possibilities: 1) Episode of Harmful Psychoactive Substance Use, 2) Harmful Pattern of Psychoactive Substance Use, and 3) Substance Dependence.
Prevention
Abuse liability
Abuse or addiction liability is the tendency to use drugs in a non-medical situation. This is typically for euphoria, mood changing, or sedation. Abuse liability is used when the person using the drugs wants something that they otherwise can not obtain. The only way to obtain this is through the use of drugs. When looking at abuse liability there are a number of determining factors in whether the drug is abused. These factors are: the chemical makeup of the drug, the effects on the brain, and the age, vulnerability, and the health (mental and physical) of the population being studied. There are a few drugs with a specific chemical makeup that leads to a high abuse liability. These are: cocaine, heroin, inhalants, marijuana, MDMA (ecstasy), methamphetamine, PCP, synthetic cannabinoids, synthetic cathinones (bath salts), nicotine (e.g. tobacco), and alcohol.
Potential vaccines for addiction to substances
Vaccines for addiction have been investigated as a possibility since the early 2000s. The general theory of a vaccine intended to "immunize" against drug addiction or other substance abuse is that it would condition the immune system to attack and consume or otherwise disable the molecules of such substances that cause a reaction in the brain, thus preventing the addict from being able to realize the effect of the drug. Addictions that have been floated as targets for such treatment include nicotine, opioids, and fentanyl. Vaccines have been identified as potentially being more effective than other anti-addiction treatments, due to "the long duration of action, the certainty of administration and a potential reduction of toxicity to important organs".
Specific addiction vaccines in development include:
NicVAX, a conjugate vaccine intended to reduce or eliminate physical dependence on nicotine. This proprietary vaccine is being developed by Nabi Biopharmaceuticals of Rockville, MD. with the support from the U.S. National Institute on Drug Abuse. NicVAX consists of the hapten 3'-aminomethylnicotine which has been conjugated (attached) to Pseudomonas aeruginosa exotoxin A.
TA-CD, an active vaccine developed by the Xenova Group which is used to negate the effects of cocaine. It is created by combining norcocaine with inactivated cholera toxin. It works in much the same way as a regular vaccine. A large protein molecule attaches to cocaine, which stimulates response from antibodies, which destroy the molecule. This also prevents the cocaine from crossing the blood–brain barrier, negating the euphoric high and rewarding effect of cocaine caused from stimulation of dopamine release in the mesolimbic reward pathway. The vaccine does not affect the user's "desire" for cocaine—only the physical effects of the drug.
TA-NIC, used to create human antibodies to destroy nicotine in the human body so that it is no longer effective.
As of September 2023, it was further reported that a vaccine "has been tested against heroin and fentanyl and is on its way to being tested against OxyContin".
Treatment
To be effective, treatment for addiction that is pharmacological or biologically based need to be accompanied by other interventions such as cognitive behavioral therapy (CBT) and dialectical behavioral therapy (DBT);individual and group psychotherapy, behavior modification strategies, twelve-step programs, and residential treatment facilities. The transtheoretical model (TTM) can be used to determine when treatment can begin and which method will be most effective. If treatment begins too early, it can cause a person to become defensive and resistant to change.
Epidemiology
Due to cultural variations, the proportion of individuals who develop a drug or behavioral addiction within a specified time period (i.e., the prevalence) varies over time, by country, and across national population demographics (e.g., by age group, socioeconomic status, etc.). Where addiction is viewed as unacceptable, there will be fewer people addicted.
Asia
The prevalence of alcohol dependence is not as high as is seen in other regions. In Asia, not only socioeconomic factors but biological factors influence drinking behavior.
Internet addiction disorder is highest in the Philippines, according to both the IAT (Internet Addiction Test) – 5% and the CIAS-R (Revised Chen Internet Addiction Scale) – 21%.
Australia
The prevalence of substance use disorder among Australians was reported at 5.1% in 2009. In 2019 the Australian Institute of Health and Welfare conducted a national drug survey that quantified drug use for various types of drugs and demographics. The survey found that in 2019, 11% of people over 14 years old smoke daily; that 9.9% of those who drink alcohol, which equates to 7.5% of the total population age 14 or older, may qualify as alcohol dependent; that 17.5% of the 2.4 million people who used cannabis in the last year may have hazardous use or a dependence problem; and that 63.5% of about 300000 recent users of meth and amphetamines were at risk for developing problem use.
Europe
In 2015, the estimated prevalence among the adult population was 18.4% for heavy episodic alcohol use (in the past 30 days); 15.2% for daily tobacco smoking; and 3.8% for cannabis use, 0.77% for amphetamine use, 0.37% for opioid use, and 0.35% for cocaine use in 2017. The mortality rates for alcohol and illicit drugs were highest in Eastern Europe. Data shows a downward trend of alcohol use among children 15 years old in most European countries between 2002 and 2014. First-time alcohol use before the age of 13 was recorded for 28% of European children in 2014.
United States
Based on representative samples of the US youth population in , the lifetime prevalence of addictions to alcohol and illicit drugs has been estimated to be approximately 8% and 2–3% respectively. Based on representative samples of the US adult population in , the 12-month prevalence of alcohol and illicit drug addictions were estimated at 12% and 2–3% respectively. The lifetime prevalence of prescription drug addictions is around 4.7%.
43.7 million people aged 12 or older surveyed by the National Survey on Drug Use and Health in the United States needed treatment for an addiction to alcohol, nicotine, or other drugs. The groups with the highest number of people were 18–25 years (25.1%) and "American Indian or Alaska Native" (28.7%). Only about 10%, or a little over 2 million, receive any form of treatments, and those that do generally do not receive evidence-based care. One-third of inpatient hospital costs and 20% of all deaths in the US every year are the result of untreated addictions and risky substance use. In spite of the massive overall economic cost to society, which is greater than the cost of diabetes and all forms of cancer combined, most doctors in the US lack the training to effectively address a drug addiction.
Estimates of lifetime prevalence rates in the US are 1–2% for compulsive gambling, 5% for sexual addiction, 2.8% for food addiction, and 5–6% for compulsive shopping. The time-invariant prevalence rate for sexual addiction and related compulsive sexual behavior (e.g., compulsive masturbation with or without pornography, compulsive cybersex, etc.) within the US ranges from 3–6% of the population.
According to a 2017 poll conducted by the Pew Research Center, almost half of US adults know a family member or close friend who has struggled with a drug addiction at some point in their life.
In 2019, opioid addiction was acknowledged as a national crisis in the United States. An article in The Washington Post stated that "America's largest drug companies flooded the country with pain pills from 2006 through 2012, even when it became apparent that they were fueling addiction and overdoses."
The National Epidemiologic Survey on Alcohol and Related Conditions found that from 2012 to 2013 the prevalence of Cannabis use disorder in U.S. adults was 2.9%.
Canada
A Statistics Canada Survey in 2012 found the lifetime prevalence and 12-month prevalence of substance use disorders were 21.6%, and 4.4% in those 15 and older. Alcohol abuse or dependence reported a lifetime prevalence of 18.1% and a 12-month prevalence of 3.2%. Cannabis abuse or dependence reported a lifetime prevalence of 6.8% and a 12-month prevalence of 3.2%. Other drug abuse or dependence has a lifetime prevalence of 4.0% and a 12-month prevalence of 0.7%. Substance use disorder is a term used interchangeably with a drug addiction.
In Ontario, Canada between 2009 and 2017, outpatient visits for mental health and addiction increased from 52.6 to 57.2 per 100 people, emergency department visits increased from 13.5 to 19.7 per 1000 people and the number of hospitalizations increased from 4.5 to 5.5 per 1000 people. Prevalence of care needed increased the most among the 14–17 age group overall.
South America
The realities of opioid use and opioid use disorder in Latin America may be deceptive if observations are limited to epidemiological findings. In the United Nations Office on Drugs and Crime report, although South America produced 3% of the world's morphine and heroin and 0.01% of its opium, prevalence of use is uneven. According to the Inter-American Commission on Drug Abuse Control, consumption of heroin is low in most Latin American countries, although Colombia is the area's largest opium producer. Mexico, because of its border with the United States, has the highest incidence of use.
Addiction and the humanities
History and etymology
The etymology of the term addiction throughout history has been misunderstood and has taken on various meanings associated with the word. An example is the usage of the word in the religious landscape of early modern Europe. "Addiction" at the time meant "to attach" to something, giving it both positive and negative connotations. The object of this attachment could be characterized as "good or bad". The meaning of addiction during the early modern period was mostly associated with positivity and goodness; during this early modern and highly religious era of Christian revivalism and Pietistic tendencies, it was seen as a way of "devoting oneself to another".
Modern research on addiction has led to a better understanding of the disease with research on the topic dating back to 1875, specifically on morphine addiction. This furthered the understanding of addiction being a medical condition. It was not until the 19th century that addiction was seen and acknowledged in the Western world as a disease, being both a physical condition and mental illness. Today, addiction is understood both as a biopsychosocial and neurological disorder that negatively impacts those who are affected by it, most commonly associated with the use of drugs and excessive use of alcohol. The understanding of addiction has changed throughout history, which has impacted and continues to impact the ways it is medically treated and diagnosed.
The suffixes "-holic" and "-holism"
In contemporary modern English "-holic" is a suffix that can be added to a subject to denote an addiction to it. It was extracted from the word alcoholism (one of the first addictions to be widely identified both medically and socially) (correctly the root "alcohol" plus the suffix "-ism") by misdividing or rebracketing it into "alco" and "-holism". There are correct medico-legal terms for such addictions: dipsomania is the medico-legal term for alcoholism; other examples are in this table:
Arts
The arts can be used in a variety of ways to address issues related to addiction. Art can be used as a form of therapy in the treatment of substance use disorders. Creative activities like painting, sculpting, music, and writing can help people express their feelings and experiences in safe and healthy ways. The arts can be used as an assessment tool to identify underlying issues that may be contributing to a person's substance use disorder. Through art, individuals can gain insights into their own motivations and behaviors that can be helpful in determining a course of treatment. Finally, the arts can be used to advocate for those suffering from a substance use disorder by raising awareness of the issue and promoting understanding and compassion. Through art, individuals can share their stories, increase awareness, and offer support and hope to those struggling with substance use disorders.
As therapy
Addiction treatment is complex and not always effective due to engagement and service availability concerns, so researchers prioritize efforts to improve treatment retention and decrease relapse rates. Characteristics of substance abuse may include feelings of isolation, a lack of confidence, communication difficulties, and a perceived lack of control. In a similar vein, people suffering from substance use disorders tend to be highly sensitive, creative, and as such, are likely able to express themselves meaningfully in creative arts such as dancing, painting, writing, music, and acting. Further evidenced by Waller and Mahony (2002) and Kaufman (1981), the creative arts therapies can be a suitable treatment option for this population especially when verbal communication is ineffective.
Primary advantages of art therapy in the treatment of addiction have been identified as:
Assess and characterize a client's substance use issues
Bypassing a client's resistances, defenses, and denial
Containing shame or anger
Facilitating the expression of suppressed and/or complicated emotions
Highlighting a client's strengths
Providing an alternative to verbal communication (via use of symbols) and conventional forms of therapy
Providing clients with a sense of control
Tackling feelings of isolation
Art therapy is an effective method of dealing with substance abuse in comprehensive treatment models. When included in psychoeducational programs, art therapy in a group setting can help clients internalize taught concepts in a more personalized manner. During the course of treatment, by examining and comparing artwork created at different times, art therapists can be helpful in identifying and diagnosing issues, as well as charting the extent or direction of improvement as a person detoxifies. Where increasing adherence to treatment regimes and maintaining abstinence is the target; art therapists can aid by customizing treatment directives (encourage the client to create collages that compare pros and cons, pictures that compare past and present and future, and drawings that depict what happened when a client went off medication).
Art therapy can function as a complementary therapy used in conjunction with more conventional therapies and can can integrate with harm reduction protocols to minimize the negative effects of drug use. An evaluation of art therapy incorporation within a pre-existing Addiction Treatment Programme based on the 12 step Minnesota Model endorsed by the Alcoholics Anonymous found that 66% of participants expressed the usefulness of art therapy as a part of treatment. Within the weekly art therapy session, clients were able to reflect and process the intense emotions and cognitions evoked by the programme. In turn, the art therapy component of the programme fostered stronger self-awareness, exploration, and externalization of repressed and unconscious emotions of clients, promoting the development of a more integrated 'authentic self'.
Despite the large number of randomized control trials, clinical control trials, and anecdotal evidence supporting the effectiveness of art therapies for use in addiction treatment, a systematic review conducted in 2018 could not find enough evidence on visual art, drama, dance and movement therapy, or 'arts in health' methodologies to confirm their effectiveness as interventions for reducing substance misuse. Music therapy was identified to have potentially strong beneficial effects in aiding contemplation and preparing those diagnosed with substance use for treatment.
As an assessment tool
The Formal Elements Art Therapy Scale (FEATS) is an assessment tool used to evaluate drawings created by people suffering from substance use disorders by comparing them to drawings of a control group (consisting of individuals without SUDs). FEATS consists of twelve elements, three of which were found to be particularly effective at distinguishing the drawings of those with SUDs from those without: Person, Realism, and Developmental. The Person element assesses the degree to which a human features are depicted realistically, the Realism element assesses the overall complexity of the artwork, and the Developmental element assesses "developmental age" of the artwork in relation to standardized drawings from children and adolescents. By using the FEATS assessment tool, clinicians can gain valuable insight into the drawings of individuals with SUDs, and can compare them to those of the control group. Formal assessments such as FEATS provide healthcare providers with a means to quantify, standardize, and communicate abstract and visceral characteristics of SUDs to provide more accurate diagnoses and informed treatment decisions.
Other artistic assessment methods include the Bird's Nest Drawing: a useful tool for visualizing a client's attachment security. This assessment method looks at the amount of color used in the drawing, with a lack of color indicating an 'insecure attachment', a factor that the client's therapist or recovery framework must take into account.
Art therapists working with children of parents suffering from alcoholism can use the Kinetic Family Drawings assessment tool to shed light on family dynamics and help children express and understand their family experiences. The KFD can be used in family sessions to allow children to share their experiences and needs with parents who may be in recovery from alcohol use disorder. Depiction of isolation of self and isolation of other family members may be an indicator of parental alcoholism.
Advocacy
Stigma can lead to feelings of shame that can prevent people with substance use disorders from seeking help and interfere with provision of harm reduction services. It can influence healthcare policy, making it difficult for these individuals to access treatment.
Artists attempt to change the societal perception of addiction from a punishable moral offense to instead a chronic illness necessitating treatment. This form of advocacy can help to relocate the fight of addiction from a judicial perspective to the public health system.
Artists who have personally lived with addiction and/or undergone recovery may use art to depict their experiences in a manner that uncovers the "human face of addiction". By bringing experiences of addiction and recovery to a personal level and breaking down the "us and them", the viewer may be more inclined to show compassion, forego stereotypes and stigma of addiction, and label addiction as a social rather than individual problem.
According to Santora the main purposes in using art as a form of advocacy in the education and prevention of substance use disorders include:
Addiction art exhibitions can come from a variety of sources, but the underlying message of these works is the same: to communicate through emotions without relying on intellectually demanding/gatekept facts and figures. These exhibitions can either stand alone, reinforce, or challenge facts.
A powerful educational tool for increasing awareness and understanding of addiction as a medical illness. Exhibitions featuring personal stories and images can help to create lasting impressions on diverse audiences (including addiction scientists/researchers, family/friends of those affected by addiction etc.), highlighting the humanity of the problem and in turn encouraging compassion and understanding.
A way to destigmatize substance use disorders and shift public perception from viewing them as a moral failing to understanding them as a chronic medical condition which requires treatment.
Provide those who are struggling with addiction assurance and encouragement of healing, and let them know that they are not alone in their struggle.
The use of visual arts can help bring attention to the lack of adequate substance use treatment, prevention, and education programs and services in a healthcare system. Messages can encourage policymakers to allocate more resources to addiction treatment and prevention from federal, state, and local levels.
The Temple University College of Public Health department conducted a project to promote awareness around opioid use and reduce associated stigma by asking students to create art pieces that were displayed on a website they created and promoted via social media. Quantitative and qualitative data was recorded to measure engagement, and the student artists were interviewed, which revealed a change in perspective and understanding, as well as greater appreciation of diverse experiences. Ultimately, the project found that art was an effective medium for empowering both the artist creating the work and the person interacting with it.
Another author critically examined works by contemporary Canadian artists that deal with addiction via the metaphor of a cultural landscape to "unmap" and "remap" ideologies related to Indigenous communities and addiction to demonstrate how colonial violence in Canada has drastically impacted the relationship between Indigenous peoples, their land, and substance abuse.
A project known as "Voice" was a collection of art, poetry and narratives created by women living with a history of addiction to explore women's understanding of harm reduction, challenge the effects of stigma and give voice to those who have historically been silenced or devalued. In the project, nurses with knowledge of mainstream systems, aesthetic knowing, feminism and substance use organized weekly gatherings, wherein women with histories of substance use and addiction worked alongside a nurse to create artistic expressions. Creations were presented at several venues, including an International Conference on Drug Related Harm, a Nursing Conference and a local gallery to positive community response.
Narrative Approach and Addiction
The narrative medicine to addiction focuses on recognizing, absorbing, and interpreting the stories of those suffering from addiction, allowing for better understanding of their experiences with narrative analysis being used to study the discourse of those with addiction. This knowledge can be used to develop better care plans with the potential to increase patient compliance and make treatment more effective.
A narrative study demonstrated and studied cognitive and emotional tendencies among substance abusers during treatment periods to learn more about motivation and ambivalence inherent in recovery over the course of a residential treatment program. Seven narrative types emerged from the overall analysis: optimistic, overly optimistic, pessimistic, overly pessimistic, "tough life," troubled/confused, and balanced. Narratives tended to express a basic level of emotionality in early stages of treatment ("optimistic", "pessimistic" narrative). Over time, as clients progressed through the program, their stories became more complex and detailed, including their drug abuse and recovery efforts, more skeptical positions towards treatment began to emerge. Clients began to distinguish between the positive and negative aspects of treatment, creating more "balanced" narratives in the process.
Due to higher medication consumption, social isolation, financial worries, and other factors, older adults are particularly vulnerable to substance use problems. Incidence of addiction among this population is inaccurately reported. Narrative therapy can provide an avenue to unearth stories of addiction in an empowering manner, and thus serves as a viable therapeutic tool in applied gerontology. When treating substance abuse in older adults, it is essential to ensure that the client is respected and comfortable disclosing information. This should be done at the outset of treatment when the therapist and older adult are developing the therapeutic relationship. The social breakdown model is an important tool that can consider the compounded effects of ageism, physical changes, social changes, and substance abuse. The narrative approach integrates the social breakdown model with substance abuse challenges and can be an effective way to address addiction in this population.
A study conducted in 2009 in the Republic of Moldova looked into the social dynamics of initiating injection drug usage by examining 42 audio-recorded, semi-structured interviews with present and former injectors. A thematic analysis suggested that self-injection was viewed as a symbolic transition of identity, enabled by interpersonal interactions and collective influences. Personal narratives of self-transition were connected to larger narratives of social transitions. The personal narratives of self-initiation and transition are contextualized and understood in terms of political (social) narratives within the core concept of the 'transitional society'. Another study examined the narratives of 'initiators': people who help people who inject drugs (PWID) with their first injection. Through their accounts, respondents described initiation events as meaningful transitions to a life characterized by predictable downfalls of homelessness, infections, and social stigma. Initiators used examples from their own personal experience to explain the process of initiation and assistance, attributing personal agency and predicting specific injection-related harms for initiates. They distinguished between two forms of harm: potentially avoidable proximal harm caused by risky injection practices (e.g. overdose, HIV) and perceived inevitable distal harm caused by long-term injection (e.g. socioeconomic decline). In this way, these narratives reflect a balance of individual agency, harm reduction intentions, and accepted notions of 'life after initiation' interact with the narrative experiences and intentions of PWIDs.
Philosophy
From a philosophy perspective, the behavior of many with addiction that is not explained by executive dysfunction or biological reasons can be explained by folk psychologyspecifically the belief–desire model. According to this model, a person acquires and uses a substance or does an addictive activity in belief that it will help them achieve a goal.
Social scientific models
Biopsychosocial–cultural–spiritual
While regarded biomedically as a neuropsychological disorder, addiction is multi-layered, with biological, psychological, social, cultural, and spiritual (biopsychosocial–cultural–spiritual) elements. A biopsychosocial–cultural–spiritual approach fosters the crossing of disciplinary boundaries, and promotes holistic considerations of addiction. A biopsychosocial–cultural–spiritual approach considers, for example, how physical environments influence experiences, habits, and patterns of addiction.
Ethnographic engagements and developments in fields of knowledge have contributed to biopsychosocial–cultural–spiritual understandings of addiction, including the work of Philippe Bourgois, whose fieldwork with street-level drug dealers in East Harlem highlights correlations between drug use and structural oppression in the United States. Prior models that have informed the prevailing biopsychosocial–cultural–spiritual consideration of addiction include:
Cultural model
The cultural model, an anthropological understanding of the emergence of drug use and abuse, was developed by Dwight Heath. Heath undertook ethnographic research and fieldwork with the Camba people of Bolivia from June 1956 to August 1957. Heath observed that adult members of society drank 'large quantities of rum and became intoxicated for several contiguous days at least twice a month'. This frequent, heavy drinking from which intoxication followed was typically undertaken socially, during festivals. Having returned in 1989, Heath observed that while much had changed, 'drinking parties' remained, as per his initial observations, and 'there appear to be no harmful consequences to anyone'. Heath's observations and interactions reflected that this form of social behavior, the habitual heavy consumption of alcohol, was encouraged and valued, enforcing social bonds in the Camba community. Despite frequent intoxication, "even to the point of unconsciousness", the Camba held no concept of alcoholism (a form of addiction), and no visible social problems associated with drunkenness, or addiction, were apparent.
As noted by Merrill Singer, Heath's findings, when considered alongside subsequent cross-cultural experiences, challenged the perception that intoxication is socially 'inherently disruptive'. Following this fieldwork, Heath proposed the 'cultural model', suggesting that 'problems' associated with heavy drinking, such as alcoholism – a recognised form addiction – were cultural: that is, that alcoholism is determined by cultural beliefs, and therefore varies among cultures. Heath's findings challenged the notion that 'continued use [of alcohol] is inexorably addictive and damaging to the consumer's health'.
The cultural model did face criticism by Sociologist Robin Room and others, who felt anthropologists could "downgrade the severity of the problem". Merrill Singer found it notable that the ethnographers working within the prominence of the cultural model were part of the 'wet generation': while not blind to the 'disruptive, dysfunctional and debilitating effects of alcohol consumption', they were products 'socialized to view alcohol consumption as normal'.
Subcultural model
Historically, addiction has been viewed from the etic perspective, defining users through the pathology of their condition. As reports of drug use rapidly increased, the cultural model found application in anthropological research exploring western drug subculture practices.
The approach evolved from the ethnographic exploration into the lived experiences and subjectivities of 1960s and 1970s drug subcultures. The seminal publication "Taking care of business", by Edward Preble and John J. Casey, documented the daily lives of New York street-based intravenous heroin users in rich detail, providing unique insight into the dynamic social worlds and activities that surrounded their drug use. These findings challenge popular narratives of immorality and deviance, conceptualizing substance abuse as a social phenomenon. The prevailing culture can have a greater influence on drug taking behaviors than the physical and psychological effects of the drug itself. To marginalized individuals, drug subcultures can provide social connection, symbolic meaning, and socially constructed purpose that they may feel is unattainable through conventional means. The subcultural model demonstrates the complexities of addiction, highlighting the need for an integrated approach. It contends that a biosocial approach is required to achieve a holistic understanding of addiction.
Critical medical anthropology model
Emerging in the early 1980s, the critical medical anthropology model was introduced, and as Merrill Singer offers 'was applied quickly to the analysis of drug use'. Where the cultural model of the 1950s looked at the social body, the critical medical anthropology model revealed the body politic, considering drug use and addiction within the context of macro level structures including larger political systems, economic inequalities, and the institutional power held over social processes.
Highly relevant to addiction, the three issues emphasized in the model are:
Self-medication
The social production of suffering
The political economy (Licit and Illicit Drugs)
These three key points highlight how drugs may come to be used to self-medicate the psychological trauma of socio-political disparity and injustice, intertwining with licit and illicit drug market politics. Social suffering, "the misery among those on the weaker end of power relations in terms of physical health, mental health and lived experience", is used by anthropologists to analyze how individuals may have personal problems caused by political and economic power. From the perspective of critical medical anthropology heavy drug use and addiction is a consequence of such larger scale unequal distributions of power.
The three models developed here – the cultural model, the subcultural model, and the Critical Medical Anthropology Model – display how addiction is not an experience to be considered only biomedically. Through consideration of addiction alongside the biological, psychological, social, cultural and spiritual (biopsychosocial–spiritual) elements which influence its experience, a holistic and comprehensive understanding can be built.
Social learning models
Social learning theory
Albert Bandura's 1977 social learning theory posits that individuals acquire addictive behaviors by observing and imitating models in their social environment. The likelihood of engaging in and sustaining similar addictive behaviors is influenced by the reinforcement and punishment observed in others. The principle of reciprocal determinism suggests that the functional relationships between personal, environmental, and behavioral factors act as determinants of addictive behavior. Thus, effective treatment targets each dynamic facet of the biopsychosocial disorder.
Transtheoretical model (stages of change model)
The transtheoretical model of change suggests that overcoming an addiction is a stepwise process that occurs through several stages.
Precontemplation: This initial stage precedes individuals considering a change in their behavior. They might be oblivious to or in denial of their addiction, failing to recognize the need for change.
Contemplation is the stage in which individuals become aware of the problems caused by their addiction and are considering change. Although they may not fully commit, they weigh the costs and benefits of making a shift.
Preparation: Individuals in this stage are getting ready to change. They might have taken preliminary steps, like gathering information or making small commitments, in preparation for behavioral change.
Action involves actively modifying behavior by making specific, observable changes to address the addictive behavior. The action stage requires significant effort and commitment.
Maintenance: After successfully implementing a change, individuals enter the maintenance stage, where they work to sustain the new behavior and prevent relapse. This stage is characterized by ongoing effort and consolidation of gains.
Termination/relapse prevention: Recognizing that relapse is a common part of the change process, this stage focuses on identifying and addressing factors that may lead to a return to old behaviors. Relapse is viewed as an opportunity for learning and strategy adjustment, with the ultimate goal of eliminating or terminating the targeted behavior.
The transtheoretical model can be helpful in guiding development of tailored behavioral interventions that can promote lasting change. Progression through these stages may not always follow a linear path, as individuals may move back and forth between stages. Resistance to change is recognized as an expected part of the process.
Addiction causes an "astoundingly high financial and human toll" on individuals and society as a whole. In the United States, the total economic cost to society is greater than that of all types of diabetes and all cancers combined. These costs arise from the direct adverse effects of drugs and associated healthcare costs (e.g., emergency medical services and outpatient and inpatient care), long-term complications (e.g., lung cancer from smoking tobacco products, liver cirrhosis and dementia from chronic alcohol consumption, and meth mouth from methamphetamine use), the loss of productivity and associated welfare costs, fatal and non-fatal accidents (e.g., traffic collisions), suicides, homicides, and incarceration, among others. The US National Institute on Drug Abuse has found that overdose deaths in the US have almost tripled among male and females from 2002 to 2017, with 72,306 overdose deaths reported in 2017 in the US. 2020 marked the year with highest number of overdose deaths over a 12-month period, with 81,000 overdose deaths, exceeding the records set in 2017.
| Biology and health sciences | Drugs and medication | null |
275712 | https://en.wikipedia.org/wiki/Cleanroom | Cleanroom | A cleanroom or clean room is an engineered space that maintains a very low concentration of airborne particulates. It is well isolated, well controlled from contamination, and actively cleansed. Such rooms are commonly needed for scientific research and in industrial production for all nanoscale processes, such as semiconductor device manufacturing. A cleanroom is designed to keep everything from dust to airborne organisms or vaporised particles away from it, and so from whatever material is being handled inside it.
A cleanroom can also prevent the escape of materials. This is often the primary aim in hazardous biology, nuclear work, pharmaceutics and virology.
Cleanrooms typically come with a cleanliness level quantified by the number of particles per cubic meter at a predetermined molecule measure. The ambient outdoor air in a typical urban area contains 35,000,000 particles for each cubic meter in the size range 0.5 μm and bigger, equivalent to an ISO 9 certified cleanroom. By comparison, an ISO 14644-1 level 1 certified cleanroom permits no particles in that size range, and just 12 particles for each cubic meter of 0.3 μm and smaller. Semiconductor facilities often get by with level 7 or 5, while level 1 facilities are exceedingly rare.
History
The modern cleanroom was invented by American physicist Willis Whitfield. As an employee of the Sandia National Laboratories, Whitfield created the initial plans for the cleanroom in 1960. Prior to Whitfield's invention, earlier cleanrooms often had problems with particles and unpredictable airflows. Whitfield designed his cleanroom with a constant, highly filtered airflow to flush out impurities. Within a few years of its invention in the 1960s, Whitfield's modern cleanroom had generated more than US$50 billion in sales worldwide (approximately $ billion today).
By mid-1963, more than 200 U.S. industrial plants had such specially constructed facilities—then using the terminology “White Rooms,” “Clean Rooms,” or “Dust-Free Rooms”—including the Radio Corporation of America, McDonnell Aircraft, Hughes Aircraft, Sperry Rand, Sylvania Electric, Western Electric, Boeing, and North American Aviation. RCA began such a conversion of part of its Cambridge, Ohio facilities in February 1961. Totalling 70,000 square feet, it was used to prepare control equipment for the Minuteman ICBM missiles.
The majority of the integrated circuit manufacturing facilities in Silicon Valley were made by three companies: MicroAire, PureAire, and Key Plastics. These competitors made laminar flow units, glove boxes, cleanrooms and air showers, along with the chemical tanks and benches used in the "wet process" building of integrated circuits. These three companies were the pioneers of the use of Teflon for airguns, chemical pumps, scrubbers, water guns, and other devices needed for the production of integrated circuits. William (Bill) C. McElroy Jr. worked as an engineering manager, drafting room supervisor, QA/QC, and designer for all three companies, and his designs added 45 original patents to the technology of the time. McElroy also wrote a four-page article for MicroContamination Journal, wet processing training manuals, and equipment manuals for wet processing and cleanrooms.
Overview
A cleanroom is a necessity in the manufacturing of semiconductors, rechargeable batteries, pharmaceutical products, and any other field that is highly sensitive to environmental contamination.
Cleanrooms can range from the very small to the very large. On the one hand, a single-user laboratory can be built to cleanroom standards within several square meters, and on the other, entire manufacturing facilities can be contained within a cleanroom with factory floors covering thousands of square meters. Between the large and the small, there are also modular cleanrooms. They have been argued to lower costs of scaling the technology, and to be less susceptible to catastrophic failure.
With such a wide area of application, not every cleanroom is the same. For example, the rooms utilized in semiconductor manufacturing need not be sterile (i.e., free of uncontrolled microbes), while the ones used in biotechnology usually must be. Vice versa, operating rooms need not be absolutely pure of nanoscale inorganic salts, such as rust, while nanotechnology absolutely requires it. What then is common to all cleanrooms is strict control of airborne particulates, possibly with secondary decontamination of air, surfaces, workers entering the room, implements, chemicals, and machinery.
Sometimes particulates exiting the compartment are also of concern, such as in research into dangerous viruses, or where radioactive materials are being handled.
Basic construction
First, outside air entering a cleanroom is filtered and cooled by several outdoor air handlers using progressively finer filters to exclude dust.
Within, air is constantly recirculated through fan units containing high-efficiency particulate absorbing filters (HEPA), and/or ultra-low particulate air (ULPA) filters to remove internally generated contaminants. Special lighting fixtures, walls, equipment and other materials are used to minimize the generation of airborne particles. Plastic sheets can be used to restrict air turbulence if the cleanroom design is of the laminar airflow type.
Air temperature and humidity levels inside a cleanroom are tightly controlled, because they affect the efficiency and means of air filtration. If a particular room requires low enough humidity to make static electricity a concern, it too will be controlled by, e.g., introducing controlled amounts of charged ions into the air using a corona discharge. Static discharge is of particular concern in the electronics industry, where it can instantly destroy components and circuitry.
Equipment inside any cleanroom is designed to generate minimal air contamination. The selection of material for the construction of a cleanroom should not generate any particulates; hence, monolithic epoxy or polyurethane floor coating is preferred. Buffed stainless steel or powder-coated mild steel sandwich partition panels and ceiling panel are used instead of iron alloys prone to rusting and then flaking. Corners like the wall to wall, wall to floor, wall to ceiling are avoided by providing coved surface, and all joints need to be sealed with epoxy sealant to avoid any deposition or generation of particles at the joints, by vibration and friction. Many cleanrooms have a "tunnel" design in which there are spaces called "service chases" that serve as air plenums carrying the air from the bottom of the room to the top so that it can be recirculated and filtered at the top of the cleanroom.
Airflow principles
Cleanrooms maintain particulate-free air through the use of either HEPA or ULPA filters employing laminar or turbulent airflow principles. Laminar, or unidirectional, airflow systems direct filtered air downward or in horizontal direction in a constant stream towards filters located on walls near the cleanroom floor or through raised perforated floor panels to be recirculated. Laminar airflow systems are typically employed across 80% of a cleanroom ceiling to maintain constant air processing. Stainless steel or other non shedding materials are used to construct laminar airflow filters and hoods to prevent excess particles entering the air. Turbulent, or non-unidirectional, airflow uses both laminar airflow hoods and nonspecific velocity filters to keep air in a cleanroom in constant motion, although not all in the same direction. The rough air seeks to trap particles that may be in the air and drive them towards the floor, where they enter filters and leave the cleanroom environment. US FDA and EU have laid down stringent guidelines and limits to ensure freedom from microbial contamination in pharmaceutical products. Plenums between air handlers and fan filter units, along with sticky mats, may also be used.
In addition to air filters, cleanrooms can also use ultraviolet light to disinfect the air. UV devices can be fitted into ceiling light fixtures and irradiate air, killing potentially infectious particulates, including 99.99 percent of airborne microbial and fungal contaminants. UV light has previously been used to clean surface contaminants in sterile environments such as hospital operating rooms. Their use in other cleanrooms may increase as equipment becomes more affordable. Potential advantages of UV-based decontamination includes a reduced reliance on chemical disinfectants and the extension of HVAC filter life.
Cleanrooms of different kinds
Some cleanrooms are kept at a positive pressure so if any leaks occur, air leaks out of the chamber instead of unfiltered air coming in. This is most typically the case in semiconductor manufacturing, where even minute amounts of particulates leaking in could contaminate the whole process, while anything leaking out . The opposite is done, e.g., in the case of high-level bio-laboratories that handle dangerous bacteria or viruses; those are always held at negative pressure, with the exhaust being passed through high-efficiency filters, and further sterilizing procedures. Both are still cleanrooms because the particulate level inside is maintained within very low limits.
Some cleanroom HVAC systems control the humidity to such low levels that extra equipment like air ionizers are required to prevent electrostatic discharge problems. This is a particular concern within the semiconductor business, because static discharge can easily damage modern circuit designs. On the other hand, active ions in the air can harm exposed components as well. Because of this, most workers in high electronics and semiconductor facilities have to wear conductive boots while working. Low-level cleanrooms may only require special shoes, with completely smooth soles that do not track in dust or dirt. However, for safety reasons, shoe soles must not create slipping hazards. Access to a cleanroom is usually restricted to those wearing a cleanroom suit, including the necessary machinery.
In cleanrooms in which the standards of air contamination are less rigorous, the entrance to the cleanroom may not have an air shower. An anteroom (known as a "gray room") is used to put on cleanroom clothing. This practice is common in many nuclear power plants, which operate as low-grade inverse pressure cleanrooms, as a whole.
Recirculating vs. one pass cleanrooms
Recirculating cleanrooms return air to the negative pressure plenum via low wall air returns. The air then is pulled by HEPA fan filter units back into the cleanroom. The air is constantly recirculating and by continuously passing through HEPA filtration removing particles from the air each time. Another advantage of this design is that air conditioning can be incorporated.
One pass cleanrooms draw air from outside and pass it through HEPA fan filter units into the cleanroom. The air then leaves through exhaust grills. The advantage of this approach is the lower cost. The disadvantages are comparatively shorter HEPA fan filter life, worse particle counts than a recirculating cleanroom, and that it cannot accommodate air conditioning.
Aseptic Practices/Processing
Aseptic practices are critical in environments where contamination control is paramount, particularly in the pharmaceutical, biotechnology, and medical device industries. Aseptic processing involves maintaining a sterile environment to prevent the introduction of contaminants during the manufacturing of products, such as sterile injectable medications and sterile medical equipment. This requires stringent control over personnel behavior, equipment sterilization, and the cleanroom environment.
There are different classifications for aseptic or sterile processing cleanrooms. The Pharmaceutical Inspection Co-operation Scheme (PIC/S) classifies cleanrooms into four grades (A, B, C, and D) based on their cleanliness level, particularly the concentration of airborne particles and viable microorganisms.
Operating procedure
In order to minimize the carrying of particulate by a person moving into the cleanroom, staff enter and leave through airlocks (sometimes including an air shower stage) and wear protective clothing such as hoods, face masks, gloves, boots, and coveralls.
Common materials such as paper, pencils, and fabrics made from natural fibers are often excluded because they shed particulates in use.
Particle levels are usually tested using a particle counter and microorganisms detected and counted through . Polymer tools used in cleanrooms must be carefully determined to be chemically compatible with cleanroom processing fluids as well as ensured to generate a low level of particle generation.
When cleaning, only special mops and buckets are used. Cleaning chemicals used tend to involve sticky elements to trap dust, and may need a second step with light molecular weight solvents to clear. Cleanroom furniture is designed to produce a minimum of particles and is easy to clean.
A cleanroom is as much a process and a meticulous culture to maintain, as it is a space as such.
Personnel contamination of cleanrooms
The greatest threat to cleanroom contamination comes from the users themselves. In the healthcare and pharmaceutical sectors, control of microorganisms is important, especially microorganisms likely to be deposited into the air stream from skin shedding. Studying cleanroom microflora is of importance for microbiologists and quality control personnel to assess changes in trends. Shifts in the types of microflora may indicate deviations from the "norm" such as resistant strains or problems with cleaning practices.
In assessing cleanroom microorganisms, the typical flora are primarily those associated with human skin (Gram-positive cocci), although microorganisms from other sources such as the environment (Gram-positive rods) and water (Gram-negative rods) are also detected, although in lower number. Common bacterial genera include Micrococcus, Staphylococcus, Corynebacterium, and Bacillus, and fungal genera include Aspergillus and Penicillium.
Cleanroom classification and standardization
Cleanrooms are classified according to the number and size of particles permitted per volume of air. Large numbers like "class 100" or "class 1000" refer to FED-STD-209E, and denote the number of particles of size 0.5 μm or larger permitted per cubic foot of air. The standard also allows interpolation; for example SNOLAB is maintained as a class 2000 cleanroom.
A discrete, light-scattering airborne particle counter is used to determine the concentration of airborne particles, equal to and larger than the specified sizes, at designated sampling locations.
Small numbers refer to ISO 14644-1 standards, which specify the decimal logarithm of the number of particles 0.1 μm or larger permitted per m3 of air. So, for example, an ISO class 5 cleanroom has at most 105 particles/m3.
Both FS 209E and ISO 14644-1 assume log-log relationships between particle size and particle concentration. For that reason, zero particle concentration does not exist. Some classes do not require testing some particle sizes, because the concentration is too low or too high to be practical to test for, but such blanks should not be read as zero.
Because 1 m3 is about 35 ft3, the two standards are mostly equivalent when measuring 0.5 μm particles, although the testing standards differ. Ordinary room air is around class 1,000,000 or ISO 9.
ISO 14644-1 and ISO 14698
ISO 14644-1 and ISO 14698 are non-governmental standards developed by the International Organization for Standardization (ISO). The former applies to cleanrooms in general (see table below), the latter to cleanrooms where biocontamination may be an issue. Since the strictest standards have been achieved only for space applications, it is sometimes difficult to know whether they were achieved in vacuum or standard conditions.
ISO 14644-1 defines the maximum concentration of particles per class and per particle size with the following formula
Where is the maximum concentration of particles in a volume of 1m of airborne particles that are equal to, or larger, than the considered particle size which is rounded to the nearest whole number, using no more than three significant figures, is the ISO class number, is the size of the particle in m and 0.1 is a constant expressed in m. The result for standard particle sizes is expressed in the following table.
US FED STD 209E
US FED-STD-209E was a United States federal standard. It was officially cancelled by the General Services Administration on November 29, 2001, but is still widely used.
Current regulating bodies include ISO, USP 800, US FED STD 209E (previous standard, still used).
Drug Quality and Security Act (DQSA) created in Nov. 2013 in response to drug compounding deaths and serious adverse events.
The Federal Food, Drug, and Cosmetic Act (FD&C Act) created specific guidelines and policies for human compounding.
503A addresses compounding by state or federally licensed facility by licensed personnel (pharmacist/ physicians)
503B pertaining to outsourcing facilities need direct supervision from a licensed pharmacist and do not need to be a licensed pharmacy. Facility is licensed through the Food and Drug Administration (FDA)
EU GMP classification
EU GMP guidelines are more stringent than others, requiring cleanrooms to meet particle counts at operation (during manufacturing process) and at rest (when manufacturing process is not carried out, but room AHU is on).
BS 5295
BS 5295 is a British Standard.
BS 5295 Class 1 also requires that the greatest particle present in any sample can not exceed 5 μm. BS 5295 has been superseded, withdrawn since the year 2007 and replaced with "BS EN ISO 14644-6:2007".
USP <800> standards
USP 800 is a United States standard developed by the United States Pharmacopeial Convention (USP) with an effective date of December 1, 2019.
Ramifications and further applications
In hospitals, theatres are similar to cleanrooms for surgical patients' operations with incisions to prevent any infections for the patient.
In another case, severely immunocompromised patients sometimes have to be held in prolonged isolation from their surroundings, for fear of infection. At the extreme, this necessitates a cleanroom environment. The same is the case for patients carrying airborne infectious diseases, only they are handled at negative, not positive pressure.
In exobiology when we seek out contact with other planets, there is a biological hazard both ways: we must not contaminate any sample return missions from other stellar bodies with terrestrial microbes, and we must not contaminate possible other ecosystems existing in other planets. Thus, even by international law, any probes we send to outer space must be sterile, and so to be handled in cleanroom conditions.
Since larger cleanrooms are very sensitive controlled environments upon which multibillion-dollar industries depend, sometimes they are even fitted with numerous seismic base isolation systems to prevent costly equipment malfunction.
| Physical sciences | Research methods | Basics and measurement |
275768 | https://en.wikipedia.org/wiki/Counting | Counting | Counting is the process of determining the number of elements of a finite set of objects; that is, determining the size of a set. The traditional way of counting consists of continually increasing a (mental or spoken) counter by a unit for every element of the set, in some order, while marking (or displacing) those elements to avoid visiting the same element more than once, until no unmarked elements are left; if the counter was set to one after the first object, the value after visiting the final object gives the desired number of elements. The related term enumeration refers to uniquely identifying the elements of a finite (combinatorial) set or infinite set by assigning a number to each element.
Counting sometimes involves numbers other than one; for example, when counting money, counting out change, "counting by twos" (2, 4, 6, 8, 10, 12, ...), or "counting by fives" (5, 10, 15, 20, 25, ...).
There is archaeological evidence suggesting that humans have been counting for at least 50,000 years. Counting was primarily used by ancient cultures to keep track of social and economic data such as the number of group members, prey animals, property, or debts (that is, accountancy). Notched bones were also found in the Border Caves in South Africa, which may suggest that the concept of counting was known to humans as far back as 44,000 BCE. The development of counting led to the development of mathematical notation, numeral systems, and writing.
Forms of counting
Verbal counting involves speaking sequential numbers aloud or mentally to track progress. Generally such counting is done with base 10 numbers: "1, 2, 3, 4", etc. Verbal counting is often used for objects that are currently present rather than for counting things over time, since following an interruption counting must resume from where it was left off, a number that has to be recorded or remembered.
Counting a small set of objects, especially over time, can be accomplished efficiently with tally marks: making a mark for each number and then counting all of the marks when done tallying. Tallying is base 1 counting.
Finger counting is convenient and common for small numbers. Children count on fingers to facilitate tallying and for performing simple mathematical operations. Older finger counting methods used the four fingers and the three bones in each finger (phalanges) to count to twelve. Other hand-gesture systems are also in use, for example the Chinese system by which one can count to 10 using only gestures of one hand. With finger binary it is possible to keep a finger count up to .
Various devices can also be used to facilitate counting, such as tally counters and abacuses.
Inclusive counting
Inclusive/exclusive counting are two different methods of counting. For exclusive counting, unit intervals are counted at the end of each interval. For inclusive counting, unit intervals are counted beginning with the start of the first interval and ending with end of the last interval. This results in a count which is always greater by one when using inclusive counting, as compared to using exclusive counting, for the same set. Apparently, the introduction of the number zero to the number line resolved this difficulty; however, inclusive counting is still useful for some things.
Refer also to the fencepost error, which is a type of off-by-one error.
Modern mathematical English language usage has introduced another difficulty, however. Because an exclusive counting is generally tacitly assumed, the term "inclusive" is generally used in reference to a set which is actually counted exclusively. For example; How many numbers are included in the set that ranges from 3 to 8, inclusive? The set is counted exclusively, once the range of the set has been made certain by the use of the word "inclusive". The answer is 6; that is 8-3+1, where the +1 range adjustment makes the adjusted exclusive count numerically equivalent to an inclusive count, even though the range of the inclusive count does not include the number eight unit interval. So, it's necessary to discern the difference in usage between the terms "inclusive counting" and "inclusive" or "inclusively", and one must recognize that it's not uncommon for the former term to be loosely used for the latter process.
Inclusive counting is usually encountered when dealing with time in Roman calendars and the Romance languages. In the ancient Roman calendar, the nones (meaning "nine") is 8 days before the ides; more generally, dates are specified as inclusively counted days up to the next named day.
In the Christian liturgical calendar, Quinquagesima (meaning 50) is 49 days before Easter Sunday. When counting "inclusively", the Sunday (the start day) will be day 1 and therefore the following Sunday will be the eighth day. For example, the French phrase for "fortnight" is quinzaine (15 [days]), and similar words are present in Greek (δεκαπενθήμερο, dekapenthímero), Spanish (quincena) and Portuguese (quinzena).
In contrast, the English word "fortnight" itself derives from "a fourteen-night", as the archaic "sennight" does from "a seven-night"; the English words are not examples of inclusive counting. In exclusive counting languages such as English, when counting eight days "from Sunday", Monday will be day 1, Tuesday day 2, and the following Monday will be the eighth day. For many years it was a standard practice in English law for the phrase "from a date" to mean "beginning on the day after that date": this practice is now deprecated because of the high risk of misunderstanding.
Similar counting is involved in East Asian age reckoning, in which newborns are considered to be 1 at birth.
Musical terminology also uses inclusive counting of intervals between notes of the standard scale: going up one note is a second interval, going up two notes is a third interval, etc., and going up seven notes is an octave.
Education and development
Learning to count is an important educational/developmental milestone in most cultures of the world. Learning to count is a child's very first step into mathematics, and constitutes the most fundamental idea of that discipline. However, some cultures in Amazonia and the Australian Outback do not count, and their languages do not have number words.
Many children at just 2 years of age have some skill in reciting the count list (that is, saying "one, two, three, ..."). They can also answer questions of ordinality for small numbers, for example, "What comes after three?". They can even be skilled at pointing to each object in a set and reciting the words one after another. This leads many parents and educators to the conclusion that the child knows how to use counting to determine the size of a set. Research suggests that it takes about a year after learning these skills for a child to understand what they mean and why the procedures are performed. In the meantime, children learn how to name cardinalities that they can subitize.
Counting in mathematics
In mathematics, the essence of counting a set and finding a result n, is that it establishes a one-to-one correspondence (or bijection) of the subject set with the subset of positive integers {1, 2, ..., n}. A fundamental fact, which can be proved by mathematical induction, is that no bijection can exist between {1, 2, ..., n} and {1, 2, ..., m} unless ; this fact (together with the fact that two bijections can be composed to give another bijection) ensures that counting the same set in different ways can never result in different numbers (unless an error is made). This is the fundamental mathematical theorem that gives counting its purpose; however you count a (finite) set, the answer is the same. In a broader context, the theorem is an example of a theorem in the mathematical field of (finite) combinatorics—hence (finite) combinatorics is sometimes referred to as "the mathematics of counting."
Many sets that arise in mathematics do not allow a bijection to be established with {1, 2, ..., n} for any natural number n; these are called infinite sets, while those sets for which such a bijection does exist (for some n) are called finite sets. Infinite sets cannot be counted in the usual sense; for one thing, the mathematical theorems which underlie this usual sense for finite sets are false for infinite sets. Furthermore, different definitions of the concepts in terms of which these theorems are stated, while equivalent for finite sets, are inequivalent in the context of infinite sets.
The notion of counting may be extended to them in the sense of establishing (the existence of) a bijection with some well-understood set. For instance, if a set can be brought into bijection with the set of all natural numbers, then it is called "countably infinite." This kind of counting differs in a fundamental way from counting of finite sets, in that adding new elements to a set does not necessarily increase its size, because the possibility of a bijection with the original set is not excluded. For instance, the set of all integers (including negative numbers) can be brought into bijection with the set of natural numbers, and even seemingly much larger sets like that of all finite sequences of rational numbers are still (only) countably infinite. Nevertheless, there are sets, such as the set of real numbers, that can be shown to be "too large" to admit a bijection with the natural numbers, and these sets are called "uncountable." Sets for which there exists a bijection between them are said to have the same cardinality, and in the most general sense counting a set can be taken to mean determining its cardinality. Beyond the cardinalities given by each of the natural numbers, there is an infinite hierarchy of infinite cardinalities, although only very few such cardinalities occur in ordinary mathematics (that is, outside set theory that explicitly studies possible cardinalities).
Counting, mostly of finite sets, has various applications in mathematics. One important principle is that if two sets X and Y have the same finite number of elements, and a function is known to be injective, then it is also surjective, and vice versa. A related fact is known as the pigeonhole principle, which states that if two sets X and Y have finite numbers of elements n and m with , then any map is not injective (so there exist two distinct elements of X that f sends to the same element of Y); this follows from the former principle, since if f were injective, then so would its restriction to a strict subset S of X with m elements, which restriction would then be surjective, contradicting the fact that for x in X outside S, f(x) cannot be in the image of the restriction. Similar counting arguments can prove the existence of certain objects without explicitly providing an example. In the case of infinite sets this can even apply in situations where it is impossible to give an example.
The domain of enumerative combinatorics deals with computing the number of elements of finite sets, without actually counting them; the latter usually being impossible because infinite families of finite sets are considered at once, such as the set of permutations of {1, 2, ..., n} for any natural number n.
| Mathematics | Basics | null |
275863 | https://en.wikipedia.org/wiki/Amborella | Amborella | Amborella is a monotypic genus of understory shrubs or small trees endemic to the main island, Grande Terre, of New Caledonia in the southwest Pacific Ocean. The genus is the only member of the family Amborellaceae and the order Amborellales and contains a single species, Amborella trichopoda. Amborella is of great interest to plant systematists because molecular phylogenetic analyses consistently place it as the sister group to all other flowering plants, meaning it was the earliest group to evolve separately from all other flowering plants.
Description
Amborella is a sprawling shrub or small tree up to high. It bears alternate, simple evergreen leaves without stipules. The leaves are two-ranked, with distinctly serrated or rippled margins, and about long.
Amborella has xylem tissue that differs from that of most other flowering plants. The xylem of Amborella contains only tracheids; vessel elements are absent. Xylem of this form has long been regarded as a primitive feature of flowering plants.
The species is dioecious. This means that each plant produces either male flowers (meaning that they have functional stamens) or female flowers (flowers with functional carpels), but not both. At any one time, a dioecious plant produces only functionally staminate or functionally carpellate flowers. Staminate ("male") Amborella flowers do not have carpels, whereas the carpellate ("female") flowers have non-functional "staminodes", structures resembling stamens in which no pollen develops. Plants may change from one reproductive morphology to the other. In one study, seven cuttings from a staminate plant produced, as expected, staminate flowers at their first flowering, but three of the seven produced carpellate flowers at their second flowering.
The small, creamy white flowers are arranged in inflorescences borne in the axils of foliage leaves. The inflorescences have been described as cymes, with up to three orders of branching, each branch being terminated by a flower. Each flower is subtended by bracts. The bracts transition into a perianth of undifferentiated tepals. The tepals typically are arranged in a spiral, but sometimes are whorled at the periphery.
Carpellate flowers are roughly in diameter, with 7 or 8 tepals. There are 1 to 3 (or rarely 0) well-differentiated staminodes and a spiral of 4 to 8 free (apocarpous) carpels. Carpels bear green ovaries; they lack a style. They contain a single ovule with the micropyle directed downwards. Staminate flowers are approximately 4 to 5 mm in diameter, with 6 to 15 tepals. These flowers bear 10 to 21 spirally arranged stamens, which become progressively smaller toward the center. The innermost may be sterile, amounting to staminodes. The stamens bear triangular anthers on short broad filaments. An anther consists of four pollen sacs, two on each side, with a small sterile central connective. The anthers have connective tips with small bumps and may be covered with secretions. These features suggest that, as with other basal angiosperms, there is a high degree of developmental plasticity.
Typically, 1 to 3 carpels per flower develop into fruit. The fruit is an ovoid red drupe (approximately 5 to 7 mm long and 5 mm wide) borne on a short (1 to 2 mm) stalk. The remains of the stigma can be seen at the tip of the fruit. The skin is papery, surrounding a thin fleshy layer containing a red juice. The inner pericarp is lignified and surrounds the single seed. The embryo is small and surrounded by copious endosperm.
Taxonomy
History
The Cronquist system, of 1981, classified the family:
Order Laurales
Subclass Magnoliidae
Class Magnoliopsida [=dicotyledons]
Division Magnoliophyta [=angiosperms]
The Thorne system (1992) classified it:
Order Magnoliales
Superorder Magnolianae
Subclass Magnoliideae [=dicotyledons]
Class Magnoliopsida [=angiosperms]
The Dahlgren system classified it:
Order Laurales
Superorder Magnolianae
Subclass Magnoliideae [=dicotyledons],
Class Magnoliopsida [=angiosperms].
Modern classification
Amborella is the only genus in the family Amborellaceae. The APG II system recognized this family, but left it unplaced at order rank due to uncertainty about its relationship to the family Nymphaeaceae. In the more recent APG systems, APG III and APG IV, the Amborellaceae comprise the monotypic order Amborellales at the base of the angiosperm phylogeny.
Phylogeny
Currently plant systematists accept Amborella trichopoda as the most basal lineage in the clade of angiosperms. In systematics the term "basal" describes a lineage that diverges near the base of a phylogeny, and thus earlier than other lineages. Since Amborella is apparently basal among the flowering plants, the features of early flowering plants can be inferred by comparing derived traits shared by the main angiosperm lineage but not present in Amborella. These traits are presumed to have evolved after the divergence of the Amborella lineage.
One early 20th century idea of "primitive" (i.e. ancestral) floral traits in angiosperms, accepted until relatively recently, is the Magnolia blossom model. This envisions flowers with numerous parts arranged in spirals on an elongated, cone-like receptacle rather than the small numbers of parts in distinct whorls of more derived flowers.
In a study designed to clarify relationships between well-studied model plants such as Arabidopsis thaliana, and the basal angiosperms Amborella, Nuphar (Nymphaeaceae), Illicium, the monocots, and more derived angiosperms (eudicots), chloroplast genomes using cDNA and expressed sequence tags for floral genes, the cladogram shown below was generated.
This hypothesized relationship of the extant seed plants places Amborella as the sister taxon to all other angiosperms, and shows the gymnosperms as a monophyletic group sister to the angiosperms. It supports the theory that Amborella branched off from the main lineage of angiosperms before the ancestors of any other living angiosperms. There is however some uncertainty about the relationship between the Amborellaceae and the Nymphaeales: one theory is that the Amborellaceae alone are the monophyletic sister to the extant angiosperms; another proposes that the Amborellaceae and Nymphaeales form a clade that is the sister group to all other extant angiosperms.
Because of its evolutionary position at the base of the flowering plant clade, there was support for sequencing the complete genome of Amborella trichopoda to serve as a reference for evolutionary studies. In 2010, the US National Science Foundation began a genome sequencing effort in Amborella, and the draft genome sequence was posted on the project website in December 2013.
Genomic and evolutionary considerations
Amborella is of great interest to plant systematists because molecular phylogenetic analyses consistently place it at or near the base of the flowering plant lineage. That is, the Amborellaceae represent a line of flowering plants that diverged very early on (more than 130 million years ago) from all the other extant species of flowering plants, and, among extant flowering plants, is the sister group to the other flowering plants. Comparing characteristics of this basal angiosperm, other flowering plants and fossils may provide clues about how flowers first appeared—what Darwin called the "abominable mystery". This position is consistent with a number of conservative characteristics of its physiology and morphology; for example, the wood of Amborella lacks the vessels characteristic of most flowering plants. The genes responsible for floral traits like scent and colors in other angiosperms, have yet to be found. Further, the female gametophyte of Amborella is even more reduced than normal female angiosperm gametophyte.
Amborella, being an understory plant in the wild, is commonly in intimate contact with shade- and moisture-dependent organisms such as algae, lichens and mosses. In those circumstances, some horizontal gene transfer between Amborella and such associated species is not surprising in principle, but the scale of such transfer has caused considerable surprise. Sequencing the Amborella mitochondrial genome revealed that for every gene of its own origin, it contains about six versions from the genomes of an assortment of the plants and algae growing with or upon it. The evolutionary and physiological significance of this is not as yet clear, nor in particular is it clear whether the horizontal gene transfer has anything to do with the apparent stability and conservatism of the species.
Ecology
Amborella is typically dioecious, but has been known to change sex in cultivation. Amborella has a mixed pollination system, relying on both insect pollinators and wind.
Conservation
The islands of New Caledonia are a biodiversity hot-spot, preserving many early diverging lineages of plants, of which Amborella is but one. This preservation has been ascribed to climate stability during and since the Tertiary (), stability that has permitted the continued survival of tropical forests on New Caledonia. In contrast, drought conditions dominated the Australian climate towards the end of the Tertiary. Current threats to biodiversity in New Caledonia include fires, mining, agriculture, invasion by introduced species, urbanization and global warming. The importance of conserving Amborella has been dramatically stated by Pillon: "The disappearance of Amborella trichopoda would imply the disappearance of a genus, a family and an entire order, as well as the only witness to at least 140 million years of evolutionary history." Conservation strategies targeted on relict species are recommended, both preserving a diversity of habitats in New Caledonia and ex situ conservation in cultivation. The IUCN conservation status is Least Concern (LC).
| Biology and health sciences | Amborellales | Plants |
275888 | https://en.wikipedia.org/wiki/Line-replaceable%20unit | Line-replaceable unit | A line-replaceable unit (LRU), lower line-replaceable unit (LLRU), line-replaceable component (LRC), or line-replaceable item (LRI) is a modular component of an airplane, ship or spacecraft (or any other manufactured device) that is designed to be replaced quickly at an operating location (1st line). The different lines (distances) are essential for logistics planning and operation. An LRU is usually a sealed unit such as a radio or other auxiliary equipment. LRUs are typically assigned logistics control numbers (LCNs) or work unit codes (WUCs) to manage logistics operations.
LRUs can improve maintenance operations, because they can be stocked and replaced quickly from distributed nearby on-site inventories (sometimes mobile storage), restoring the mobile systems to service, while the failed (unserviceable) LRU is undergoing complicated repair and overhaul actions in other support locations (lines). Because of their modularity, LRUs also can contribute reducing system costs and increase quality, by centralizing development across different models of vehicles.
LRUs are similar in nature to shop-replaceable units (SRUs), but rather than being component functions, represent complete functional units.
Definition
While the term LRU has been in use for decades, MIL-PRF-49506, Notice 1 of 18 Jan 05, the Performance Spec for Logistics Management Information defines an LRU as:
An LRU is an essential support item which is removed and replaced at the field level to restore the end item to an operational ready condition. Conversely, a non-LRU is a part, component, or assembly used in the repair of an LRU / LLRU, when the LRU has failed and has been removed from the end item for repair.
An LLRU is part of an LRU, and which can be removed and replaced at the field level to restore its LRU to an operational ready condition. As an LRU is considered the 'parent', the LLRU is considered a 'child'. An LLRU can also be a child of a child--that is, an LLRU being a component of a higher-level LLRU. However, there is no hierarchy difference between child levels; the only hierarchical separation is parent versus child.
Specifications
LRUs are designed to specifications to assure that they can be interchanged, especially if they are from different manufacturers. Usually a class of LRUs will have coordinated environmental specifications (i.e. temperature, condensation, etc.). However, each particular LRU will also have detailed specifications describing its function, tray size, tray connectors, attachment points, weight ranges, etc. It is common for LRU trays to have standardized connections for rapid mounting, cooling air, power, and grounding. The mounting hardware is often manually removable standard-screw-detent quick-release fittings. Front-mounted electrical connectors are often jacks for ring-locked cannon plugs that can be removed and replaced (R&R) without tools. Specifications also define the supporting tools necessary to remove and replace the unit. Many require no tools, or a standard-sized Frearson screwdriver. Frearson is specified for some vehicles and many marine systems because Frearson screws keep their mating screwdriver from camming out, and the same screwdriver can be used on many sizes of screws. Most LRUs also have handles, and specific requirements for their bulk and weight. LRUs typically need to be "transportable" and fit through a door or hatchway. There are also requirements for flammability, unwanted radio emissions, resistance to damage from fungus, static electricity, heat, pressure, humidity, condensation drips, vibration, radiation, and other environmental measurements.
LRUs may be designed to ARINC 700-series standards. The form factor of LRUs comply to ARINC Standards, ARINC 404 and ARINC 600. LRUs are also defined by manufacturers like Airbus and Boeing and by various military organizations. In the military, electronic LRUs are typically designed to interface according to data bus standards such as MIL-STD-1553. On the International Space Station, LRUs are referred to as Orbit Replaceable Units.
| Technology | Aircraft components | null |
275978 | https://en.wikipedia.org/wiki/Camphor | Camphor | Camphor () is a waxy, colorless solid with a strong aroma. It is classified as a terpenoid and a cyclic ketone. It is found in the wood of the camphor laurel (Cinnamomum camphora), a large evergreen tree found in East Asia; and in the kapur tree (Dryobalanops sp.), a tall timber tree from South East Asia. It also occurs in some other related trees in the laurel family, notably Ocotea usambarensis. Rosemary leaves (Rosmarinus officinalis) contain 0.05 to 0.5% camphor, while camphorweed (Heterotheca) contains some 5%. A major source of camphor in Asia is camphor basil (the parent of African blue basil). Camphor can also be synthetically produced from oil of turpentine.
The compound is chiral, existing in two possible enantiomers as shown in the structural diagrams. The structure on the left is the naturally occurring (+)-camphor ((1R,4R)-bornan-2-one), while its mirror image shown on the right is the (−)-camphor ((1S,4S)-bornan-2-one). Camphor has few uses but is of historic significance as a compound that is readily purified from natural sources.
Etymology
The word camphor derived in the 14th century from Old , itself from Medieval , from , perhaps through , from apparently from Austronesian 'lime' (chalk).
In Old Malay, camphor was called , meaning "the chalk of Barus", referring to Barus, an ancient port near modern Sibolga on the western coast of Sumatra. This port traded in camphor extracted from the Borneo camphor trees (Dryobalanops aromatica) that were abundant in the region.
Production
Natural camphor
(+)-camphor has been produced as a forest product for centuries, condensed from the vapor given off by the roasting of wood chips cut from Camphora officinarum, and later by passing steam through the pulverized wood and condensing the vapors. By the early 19th century most camphor tree reserves had been depleted with the remaining large stands in Japan and Taiwan, with Taiwanese production greatly exceeding Japanese. Camphor was one of the primary resources extracted by Taiwan's colonial powers as well as one of the most lucrative. First the Chinese and then the Japanese established monopolies on Taiwanese camphor. In 1868, a British naval force sailed into Anping harbor and the local British representative demanded the end of the Chinese camphor monopoly. After the local imperial representative refused, the British bombarded the town and took the harbor. The "camphor regulations" negotiated between the two sides subsequently saw a brief end to the camphor monopoly.
(-)-camphor occurs naturally in the essential oil of Matricaria plants. As a result, it's much rarer.
Synthetic camphor
Camphor is produced from alpha-pinene, which is abundant in the oils of coniferous trees and can be distilled from turpentine produced as a side product of chemical pulping. With acetic anhydride as the solvent and with catalysis by a strong acid, alpha-pinene is converted to isobornyl acetate. Hydrolysis of this ester gives isoborneol which can be oxidized to give racemic camphor.
A biological enzyme has been proposed for producing only the rare (-) or L-camphor. This EstB esterase from Burkholderia gladioli hydrolyzes only (+)-isobornyl acetate.
Reactions
The reactions of camphor have been extensively examined. Some representative transformations include
sulfonation:
oxidation with selenium dioxide to camphorquinone .
Camphor can also be reduced to isoborneol using sodium borohydride.
Biochemistry
Biosynthesis
In biosynthesis, camphor is produced from geranyl pyrophosphate, via cyclisation of linaloyl pyrophosphate to bornyl pyrophosphate, followed by hydrolysis to borneol and oxidation to camphor.
Uses
The first significant manmade plastics were low-nitrogen (or "soluble") nitrocellulose (pyroxylin) plastics. In the early decades of the plastics industry, camphor was used in immense quantities as the plasticizer that creates celluloid from nitrocellulose, in nitrocellulose lacquers and other plastics and lacquers.
Alternative medicine and scent
Camphor has been used for its scent, as an embalming fluid, as topical medication, as a manufacturing chemical, and in religious ceremonies.
Camphor has been used as a folk medicine over centuries, probably most commonly as a decongestant. Camphor was used in ancient Sumatra to treat sprains, swellings, and inflammation. Camphor also was used for centuries in traditional Chinese medicine for various purposes. In Europe, camphor was used after the Black Death era.
In the 20th century, camphor was used as an analeptic by injection, and to induce seizures in schizophrenic people in an attempt to treat psychosis.
Camphor has limited use in veterinary medicine by intramuscular injection to treat breathing difficulties in horses.
Topical medication
Camphor is commonly applied as a topical medication as a skin cream or ointment to relieve itching from insect bites, minor skin irritation, or joint pain. It is absorbed in the skin epidermis, where it stimulates nerve endings sensitive to heat and cold, producing a warm sensation when vigorously applied, or a cool sensation when applied gently, indicating its properties as a counterirritant. The action on nerve endings also induces a slight local analgesia.
Respiratory aerosol
Camphor is also used via an aerosol, typically by steam inhalation, sometimes in the form of branded nasal inhaler sticks, to inhibit coughing and relieve upper airway congestion due to the common cold. However, the clinical efficacy of these remedies is challenged.
Other niche uses
Camphor is used by marksmen to blacken the front and rear sights of rifles to prevent the sights from reflecting. This is done by setting light to a small amount of camphor, which burns at a relatively low temperature, and using the soot rising from the flame to deposit a coating on a surface held above it. Historically, this soot blackening was also used to coat barograph record charts. A barracks-room rumour possibly derived from an older rumour about saltpetre holds that soldiers in the Islamic Republic of Iran are dosed with camphor in their daily meals in order to repress their libido and prevent homosexual incidents.
Pest deterrent and preservative
Camphor is believed to be toxic to insects and is thus sometimes used as a repellent. Camphor is used as an alternative to mothballs. Camphor crystals are sometimes used to prevent damage to insect collections by other small insects. It is kept in clothes used on special occasions and festivals, and also in cupboard corners as a cockroach repellent. The smoke of camphor crystal or camphor incense sticks can be used as an environmentally-friendly mosquito repellent.
Recent studies have indicated that camphor essential oil can be used as an effective fumigant against red fire ants, as it affects the attacking, climbing, and feeding behavior of major and minor workers.
Camphor is also used as an antimicrobial substance. In embalming, camphor oil was one of the ingredients used by ancient Egyptians for mummification.
Perfume
In the ancient Arab world, camphor was a common perfume ingredient. The Chinese referred to the best camphor as "dragon's brain perfume", due to its "pungent and portentous aroma" and "centuries of uncertainty over its provenance and mode of origin".
Culinary uses
One of the earliest known recipes for ice cream dating to the Tang dynasty includes camphor as an ingredient. It was used to flavor leavened bread in ancient Egypt. In ancient and medieval Europe, camphor was used as an ingredient in sweets. It was used in a wide variety of both savory and sweet dishes in medieval Arabic language cookbooks, such as Kitab al-Ṭabikh compiled by ibn Sayyār al-Warrāq in the 10th century. It also was used in sweet and savory dishes in the Ni'matnama, according to a book written in the late 15th century for the sultans of Mandu. It is a main constituent of a spice known as "edible camphor" (or kapur), which may be used in traditional South Indian desserts like Payasam and Chakkarai Pongal.
Religious rites
Camphor is widely used in Hindu religious ceremonies. Aarti is performed after placing it on a stand and setting fire to it usually as the last step of puja or devotional worship ritual to one or more deities. Camphor is mentioned in the Quran as being the fragrance of wine given to believers in heaven.
Toxicity
Applied on skin, camphor may cause allergic reactions in some people; when ingested by mouth, camphor cream or ointment is poisonous. In high ingested doses, camphor produces symptoms of irritability, disorientation, lethargy, muscle spasms, vomiting, abdominal cramps, convulsions, and seizures. Lethal doses by ingestion in adults are in the range 50–500 mg/kg (orally). Generally, ingestion of two grams causes serious toxicity and four grams is potentially lethal.
Airborne camphor may be toxic if respired by humans. The Permissible Exposure Limit (PEL) for camphor in ambient air is 2 mg/m3 at exposure time (TWA) not more than 8 hours. 200 mg/m3 is considered a very dangerous concentration (IDLH).
History of synthetic camphor
When its use in the nascent chemical industries (discussed below) greatly increased the volume of demand in the late 19th century, potential for changes in supply and in price followed. In 1911 Robert Kennedy Duncan, an industrial chemist and educator, related that the Imperial Japanese government had recently (1907–1908) tried to monopolize the production of natural camphor as a forest product in Asia but that the monopoly was prevented by the development of the total synthesis alternatives, which began in "purely academic and wholly uncommercial" form with Gustav Komppa's first report:
This ongoing check on price growth was confirmed in 1942 in a monograph on DuPont's history, where William S. Dutton said, "Indispensable in the manufacture of pyroxylin plastics, natural camphor imported from Formosa and selling normally for about 50 cents a pound, reached the high price of $3.75 in 1918 [amid the global trade disruption and high explosives demand that World War I created]. The organic chemists at DuPont replied by synthesizing camphor from the turpentine of southern US pine stumps, with the result that the price of industrial camphor sold in carload lots in 1939 was between 32 cents and 35 cents a pound."
The background of Gustaf Komppa's synthesis was as follows. In the 19th century, it was known that nitric acid oxidizes camphor into camphoric acid. Haller and Blanc published a semisynthesis of camphor from camphoric acid. Although they demonstrated its structure, they were unable to prove it. The first complete total synthesis of camphoric acid was published by Komppa in 1903. Its inputs were diethyl oxalate and 3,3-dimethylpentanoic acid, which reacted by Claisen condensation to yield diketocamphoric acid. Methylation with methyl iodide and a complicated reduction procedure produced camphoric acid. William Perkin published another synthesis a short time later. Previously, some organic compounds (such as urea) had been synthesized in the laboratory as a proof of concept, but camphor was a scarce natural product with a worldwide demand. Komppa realized this, and began industrial production of camphor in Tainionkoski, Finland, in 1907 (with plenty of competition, as Kennedy Duncan reported).
A different way of synthesis was developed at the same time by Dr. Karl Stephan from Chemische Fabrik auf Actien. This chemist, who had patented a route to synthesize camphene in 1902, found out that borneol or isoborneol could easily be oxidized with permanganate in benzene solution with unprecedentedly high yields of 95+%, and patented it in 1903. The process was efficient enough to compete with natural camphor, and Japan was forced to lower prices in 1907, but the German company still increased its production, reaching 623 tons in 1913, only to be interrupted by WWI.
| Physical sciences | Terpenes and terpenoids | Chemistry |
275991 | https://en.wikipedia.org/wiki/Range%20of%20a%20function | Range of a function | In mathematics, the range of a function may refer to either of two closely related concepts:
the codomain of the function, or
the image of the function.
In some cases the codomain and the image of a function are the same set; such a function is called surjective or onto. For any non-surjective function the codomain and the image are different; however, a new function can be defined with the original function's image as its codomain, where This new function is surjective.
Definitions
Given two sets and , a binary relation between and is a function (from to ) if for every element in there is exactly one in such that relates to . The sets and are called the domain and codomain of , respectively. The image of the function is the subset of consisting of only those elements of such that there is at least one in with .
Usage
As the term "range" can have different meanings, it is considered a good practice to define it the first time it is used in a textbook or article. Older books, when they use the word "range", tend to use it to mean what is now called the codomain. More modern books, if they use the word "range" at all, generally use it to mean what is now called the image. To avoid any confusion, a number of modern books don't use the word "range" at all.
Elaboration and example
Given a function
with domain , the range of , sometimes denoted or , may refer to the codomain or target set (i.e., the set into which all of the output of is constrained to fall), or to , the image of the domain of under (i.e., the subset of consisting of all actual outputs of ). The image of a function is always a subset of the codomain of the function.
As an example of the two different usages, consider the function as it is used in real analysis (that is, as a function that inputs a real number and outputs its square). In this case, its codomain is the set of real numbers , but its image is the set of non-negative real numbers , since is never negative if is real. For this function, if we use "range" to mean codomain, it refers to ; if we use "range" to mean image, it refers to .
For some functions, the image and the codomain coincide; these functions are called surjective or onto. For example, consider the function which inputs a real number and outputs its double. For this function, both the codomain and the image are the set of all real numbers, so the word range is unambiguous.
Even in cases where the image and codomain of a function are different, a new function can be uniquely defined with its codomain as the image of the original function. For example, as a function from the integers to the integers, the doubling function is not surjective because only the even integers are part of the image. However, a new function whose domain is the integers and whose codomain is the even integers is surjective. For the word range is unambiguous.
| Mathematics | Functions: General | null |
275992 | https://en.wikipedia.org/wiki/Caracal | Caracal | The caracal (Caracal caracal) () is a medium-sized wild cat native to Africa, the Middle East, Central Asia, and arid areas of Pakistan and northwestern India. It is characterised by a robust build, long legs, a short face, long tufted ears, relatively short tail, and long canine teeth. Its coat is uniformly reddish tan or sandy, while the ventral parts are lighter with small reddish markings. It reaches at the shoulder and weighs . It was first scientifically described by German naturalist Johann Christian Daniel von Schreber in 1776. Three subspecies are recognised.
Typically nocturnal, the caracal is highly secretive and difficult to observe. It is territorial, and lives mainly alone or in pairs. The caracal is a carnivore that typically preys upon birds, rodents, and other small mammals. It can leap higher than and catch birds in midair. It stalks its prey until it is within of it, after which it runs it down and kills it with a bite to the throat or to the back of the neck. Both sexes become sexually mature by the time they are one year old and breed throughout the year. Gestation lasts between two and three months, resulting in a litter of one to six kittens. Juveniles leave their mothers at the age of nine to ten months, though a few females stay back with their mothers. The average lifespan of captive caracals is nearly 16 years.
Etymology
The name 'caracal' was proposed by Georges Buffon in 1761 who referred to its Turkish name 'Karrah-kulak' or 'Kara-coulac', meaning 'black ear'. The 'lynx' of the Greeks and Romans was most probably the caracal, and the name 'lynx' is sometimes still applied to it, but the present-day lynx proper is a separate genus.
The caracal is also known as desert lynx and Persian lynx.
Taxonomy and phylogeny
Felis caracal was the scientific name used by Johann Christian Daniel von Schreber in 1776 who described a caracal skin from the Cape of Good Hope. In 1843, John Edward Gray placed it in the genus Caracal. It is placed in the family Felidae and subfamily Felinae.
In the 19th and 20th centuries, several caracal specimens were described and proposed as subspecies. Since 2017, three subspecies have been recognised as valid:
Southern caracal (C. c. caracal) – occurs in Southern and East Africa
Northern caracal (C. c. nubicus) – occurs in North and West Africa
Asiatic caracal (C. c. schmitzi) – occurs in Asia
Phylogeny
Results of a phylogenetic study indicates that the caracal and the African golden cat (Caracal aurata) diverged between 2.93 and 1.19 million years ago. These two species together with the serval (Leptailurus serval) form the Caracal lineage, which diverged between 11.56 and 6.66 million years ago. The ancestor of this lineage arrived in Africa between 8.5 and 5.6 million years ago.
The relationship of the caracal is considered as follows:
Characteristics
The caracal is a slender, moderately sized cat characterised by a robust build, a short face, long canine teeth, tufted ears, and long legs. It reaches nearly at the shoulder. The tan, bushy tail extends to the hocks. The caracal is sexually dimorphic; the females are smaller than the males in most bodily parameters.
The prominent facial features include the long black tufts on the ears, two black stripes from the forehead to the nose, the black outline of the mouth, the distinctive black facial markings, and the white patches surrounding the eyes and the mouth. The eyes appear to be narrowly open due to the lowered upper eyelid, probably an adaptation to shield the eyes from the sun's glare. The ear tufts may start drooping as the animal ages. The coat is uniformly reddish tan or sandy, though black caracals are also known. The underbelly and the insides of the legs are lighter, often with small reddish markings. The fur, soft, short, and dense, grows coarser in the summer. The ground hairs are denser in winter than in summer. The length of the guard hairs can be up to long in winter, but shorten to in summer. These features indicate the onset of moulting in the hot season, typically in October and November. The hind legs are longer than the forelegs, so the body appears to be sloping downward from the rump.
Male caracals measure in head-to-body length and have long tails; 77 male caracals ranged in weight between . The head-to-body length of females is with a tail of ; 63 females ranged in weight between .
The caracal is often confused with a lynx, as both cats have tufted ears. However, a notable point of difference between the two is that Lynx species are spotted and blotched, while the caracal shows no such markings on the coat. The African golden cat has a similar build as the caracal's, but is darker and lacks the ear tufts. The sympatric serval can be distinguished from the caracal by the former's lack of ear tufts, white spots behind the ears, spotted coat, longer legs, longer tail, and smaller footprints.
The skull of the caracal is high and rounded, featuring large auditory bullae, a well-developed supraoccipital crest normal to the sagittal crest, and a strong lower jaw. The caracal has a total of 30 teeth; the dental formula is . The deciduous dentition is . The canines are up to long and sharp. The caracal lacks the second upper premolars, and the upper molars are diminutive. The large paws have four digits in the hind legs and five in the fore legs. The first digit of the fore leg remains above the ground and features the dewclaw. The sharp and retractile claws are larger but less curved in the hind legs.
Distribution and habitat
In Africa, the caracal is widely distributed south of the Sahara, but considered rare in North Africa. In Asia, it occurs from the Arabian Peninsula, Middle East, Turkmenistan, Uzbekistan to western India. It inhabits forests, savannas, marshy lowlands, semideserts, and scrub forests, but prefers dry areas with low rainfall and availability of cover. In montane habitats such as in the Ethiopian Highlands, it occurs up to an elevation of .
In Ethiopia's Degua Tembien massif, they can be seen along roads, sometimes as roadkills.
In the Emirate of Abu Dhabi, a male caracal was photographed by camera traps in Jebel Hafeet National Park in Al Ain Region, in spring 2019, the first such record since the mid-1980s.
In south-western Turkey, caracals are present in a wildlife reserve in the provinces of Antalya and Muğla that is regularly used for recreational and agricultural activities. However, they avoid humans and are active mostly at night. In Uzbekistan, caracals were recorded only in the desert regions of the Ustyurt Plateau and Kyzylkum Desert. Between 2000 and 2017, 15 individuals were sighted alive, and at least 11 were killed by herders.
In Iran, the caracal has been recorded in Abbasabad Naein Reserve, Bahram’gur Protected Area, Kavir National Park and in Yazd province.
In India, the caracal occurs in Sariska Tiger Reserve and Ranthambhore National Park. The Indian population may be under 100, and is thought extinct in 10 of the 13 Indian states it had historical populations in.
Ecology and behaviour
The caracal is typically nocturnal, though some activity may be observed during the day as well. However, the cat is so secretive and difficult to observe that its activity at daytime might easily go unnoticed. A study in South Africa showed that caracals are most active when the air temperature drops below ; activity typically ceases at higher temperatures. A solitary cat, the caracal mainly occurs alone or in pairs; the only groups seen are of mothers with their offspring. Females in oestrus temporarily pair with males. A territorial animal, the caracal marks rocks and vegetation in its territory with urine and probably with dung, which is not covered with soil. Claw scratching is prominent, and dung middens are typically not formed. In Israel, males are found to have territories averaging , while that of females averaged . The male territories vary from in Saudi Arabia. In Mountain Zebra National Park, the home ranges of females vary between . These territories overlap extensively. The conspicuous ear tufts and the facial markings often serve as a method of visual communication; caracals have been observed interacting with each other by moving the head from side to side so that the tufts flicker rapidly. Like other cats, the caracal meows, growls, hisses, spits, and purrs.
Diet and hunting
The caracal typically preys on mammals, which generally comprise at least 80% of its diet; and it also targets small Cercopithecidae monkeys and larger antelopes such as young kudu, impala, mountain gazelle, dorcas gazelle, Cape bushbuck, gerenuk, mountain reedbuck, Sharpe's grysbok and springbok. The remaining percentage is made up of lizards, snakes and insects. Rodents comprise a significant portion of its diet in western India. It tends to focus on the most abundant prey species. Occasionally, it consumes grasses and grapes, which help to clear the immune system and stomach of any parasites.
In South Africa, caracals prey on Cape grysbok, common duiker, bush vlei rats, rock hyrax and Cape hare. In areas where sheep and goat are farmed, caracals have also been documented preying on small livestock, but this is however only a small and seasonal portion of their diet when wild prey is scarce. Caracals are estimated to have caused about 11% of African penguin mortality in Simon's Town breeding colony between January 1999 and September 2021.
The caracals's speed and agility make it an efficient hunter, able to take down prey two to three times its size. The powerful hind legs allow it to leap more than in the air to catch birds on the wing. It can even twist and change its direction mid-air. It is an adroit climber. It stalks its prey until it is within , following which it can launch into a sprint. While large prey such as antelopes are suffocated by a throat bite, smaller prey are killed by a bite on the back of the neck. Kills are consumed immediately, and less commonly dragged to cover. It returns to large kills if undisturbed. It has been observed to begin feeding on antelope kills at the hind parts. It may scavenge at times, though this has not been frequently observed.
Reproduction
Both sexes become sexually mature by the time they are a year old; production of gametes begins even earlier at seven to ten months. However, successful mating takes place only at 12 to 15 months. Breeding takes place throughout the year. Oestrus, one to three days long, recurs every two weeks unless the female is pregnant. Females in oestrus show a spike in urine-marking, and form temporary pairs with males. Mating has not been extensively studied; a limited number of observations suggest that copulation, lasting nearly four minutes on an average, begins with the male smelling the areas urine-marked by the female, which then rolls on the ground. He then approaches and mounts her. The pair separate after copulation.
Gestation lasts about two to three months, following which a litter consisting of one to six kittens is born. Births generally peak from October to February. Births take place in dense vegetation or deserted burrows of aardvarks and porcupines. Kittens are born with their eyes and ears shut and the claws not retractable (unable to be drawn inside); the coat resembles that of adults, but the abdomen is spotted. Eyes open by ten days, but it takes longer for the vision to become normal. The ears become erect and the claws become retractable by the third or the fourth week. Around the same time, the kittens start roaming their birthplace, and start playing among themselves by the fifth or the sixth week. They begin taking solid food around the same time; they have to wait for nearly three months before they make their first kill. As the kittens start moving about by themselves, the mother starts shifting them every day. All the milk teeth appear in 50 days, and permanent dentition is completed in 10 months. Juveniles begin dispersing at nine to ten months, though a few females stay back with their mothers. The average lifespan of the caracal in captivity is nearly 16 years.
In the 1990s, a captive caracal spontaneously mated with a domestic cat in the Moscow Zoo, resulting in a felid hybrid offspring.
Threats
The caracal is listed as Least Concern on the IUCN Red List since 2002, as it is widely distributed in over 50 range countries, where the threats to caracal populations vary in extent. Habitat loss due to agricultural expansion, the building of roads and settlements is a major threat in all range countries. It is thought to be close to extinction in North Africa, critically endangered in Pakistan, endangered in Jordan, but stable in central and Southern Africa. Local people kill caracal to protect livestock, or in retaliation for its preying on small livestock. Additionally, it is threatened by hunting for the pet trade on the Arabian Peninsula. In Turkey and Iran, caracals are frequently killed in road accidents.
In Uzbekistan, the major threat to caracal is killing by herders in retaliation for livestock losses. Guarding techniques and sheds are inadequate to protect small livestock like goats and sheep from being attacked by predators. Additionally, similarly to Ethiopia, heavy-traffic roads crossing caracal habitat pose a potential threat for the species.
Conservation
African caracal populations are listed under CITES Appendix II, while Asian populations come under CITES Appendix I. Hunting of caracal is prohibited in Afghanistan, Algeria, Egypt, India, Iran, Israel, Jordan, Kazakhstan, Lebanon, Morocco, Pakistan, Syria, Tajikistan, Tunisia, Turkey, Turkmenistan, and Uzbekistan. Caracals occur in a number of protected areas across their range.
In South Africa, the caracal is considered Least Concern, as it is widespread and adaptable to a variety of habitats. It is tolerant to human-dominated areas, although it has been persecuted for many decades. Farmers are encouraged to report sightings of caracals, both dead and alive, and livestock killed by caracals to the national Predation Management Information Centre.
The Central Asian caracal population is listed as Critically Endangered in Uzbekistan since 2009, and in Kazakhstan since 2010.
In culture
The caracal appears to have been religiously significant in the ancient Egyptian culture, as it occurs in paintings and as bronze figurines; sculptures are thought to have guarded the tombs of pharaohs. Embalmed caracals have also been discovered.
The caracal was esteemed for its ability to catch birds in flight and was used for coursing by Mughal emperors in India at least since the Delhi Sultanate. Chinese emperors used caracals as gifts. In the 13th and the 14th centuries, Yuan dynasty rulers bought numerous caracals, cheetahs (Acinonyx jubatus) and tigers (Panthera tigris) from Muslim merchants in the western parts of the empire in return for gold, silver, copper cash and silk. According to the Ming Shilu, the subsequent Ming dynasty continued this practice. Until the 20th century, the caracal was used in hunts of Indian rulers to hunt small game, while the cheetah was used for larger game. In those times, caracals were used to hunt bustards, francolins, and other game birds. They were also placed in arenas with flocks of pigeons and people would bet on which caracal would kill the largest number of pigeons. This probably gave rise to the expression "to put the cat among the pigeons". Its pelt was used for making fur coats.
| Biology and health sciences | Felines | Animals |
276041 | https://en.wikipedia.org/wiki/Palomino | Palomino | Palomino is a genetic color in horses, consisting of a gold coat and white mane and tail; the degree of whiteness can vary from bright white to yellow. The palomino color derived from the inter-breeding of Spanish horses with those from the United States. Genetically, the palomino color is created by a single allele of a dilution gene called the cream gene working on a "red" (chestnut) base coat. Palomino is created by a genetic mechanism of incomplete dominance, hence it is not considered true-breeding. However, most color breed registries that record palomino horses were founded before equine coat color genetics were understood as well as they are today, therefore the standard definition of a palomino is based on the visible coat color, not heritability nor the underlying presence of the dilution gene.
Due to their distinct color, palominos stand out in a show ring, and are much sought after as parade horses. They were particularly popular in movies and television during the 1940s and 1950s. One of the most famous palomino horses was Trigger, known as "the smartest horse in movies", the faithful mount of the Hollywood cowboy star Roy Rogers. Another famous palomino was Mister Ed (real name Bamboo Harvester) who starred on his own TV show in the 1960s. A palomino was also featured in the show Xena: Warrior Princess (1995–2001). Xena's horse Argo was portrayed by a palomino mare named Tilly. In today's horse breeding the palomino color can be created by crossing a chestnut with a cremello.
Palomino is a Spanish word meaning juvenile pigeon (the diminutive of paloma, pigeon) and its equine usage refers to the color of such birds.
Description
Palomino horses have a yellow or gold coat, with a white or light cream mane and tail. The shades of the body coat color range from cream to a dark gold.
Unless also affected by other, unrelated genes, palominos have dark skin and brown eyes, though some may be born with pinkish skin that darkens with age. Some have slightly lighter brown or amber eyes. A heterozygous cream dilute (CR) such as the palomino must not be confused with a horse carrying champagne dilution. Champagne (CH) dilutes are born with pumpkin-pink skin and blue eyes, which darken within days to amber, green or light brown, and their skin acquires a darker mottled complexion around the eyes, muzzle, and genitalia as the animal matures.
A horse with rosy-pink skin and blue eyes in adulthood is most often a cremello or a perlino, a horse carrying two cream dilution genes.
Sooty palomino horses may have darker hairs in the mane, tail and coat. The summer coat of a palomino is usually a slightly darker shade than the winter coat.
Colors confused with palomino
Many non-palominos may also have a gold or tan coat and a light mane and tail.
Chestnut with flaxen mane and tail: Lighter chestnuts with a light cream mane and tail carry a flaxen gene, but not a cream dilution. For example, the Haflinger breed has many light chestnuts with flaxen that may superficially resemble dark palomino, but there is no cream gene in the breed.
Cremellos carry two copies of the cream gene and have a light mane and tail but also a cream-colored hair coat, rosy pink skin and blue eyes.
The champagne gene is the most similar palomino mimic, as it creates a golden-colored coat on some horses, but golden champagnes have light skin with mottling, blue eyes at birth, and amber or hazel eyes in adulthood.
Horses with a very dark brown coat but a flaxen mane and tail are sometimes called "chocolate palomino", and some palomino color registries accept horses of such color. However, this coloring is not genetically palomino. There are two primary ways the color is created. The best-known is a liver chestnut with a flaxen mane and tail. The genetics that create light flaxen manes and tails on otherwise chestnut horses are not yet fully understood, but they are not the same as the cream dilution. The other genetic mechanism is derived from the silver dapple gene, which lightens a black coat to dark brown, and affects the mane and tail even more strongly, diluting to cream or near-white.
Buckskins have a golden body coat but a black mane and tail. Buckskin is also created by the action of a single cream gene, but on a bay coat.
Dun horses have a tan body with a darker mane and tail plus primitive markings such as a dorsal stripe down the spine and horizontal striping on the upper back of the forearm.
The pearl gene in a homozygous state creates a somewhat apricot-colored coat with pale skin. When crossed with a single cream gene, the resulting horse, often called a "pseudo-double-dilute", appears visually to be a cremello.
Color breed registries
In the United States, some palomino horses are classified as a color breed. However, unlike the Appaloosa or the Friesian, which are distinct breeds that also happen to have a unique color preference, Palomino color breed registries often accept a wide range of breed or type if the animals are properly golden-colored. The Palomino cannot be a true horse breed, however, because palomino color is an incomplete dominant gene and does not breed "true". A palomino crossed with a palomino may result in a palomino about 50% of the time, but could also produce a chestnut (25% probability) or a cremello (25% probability). Thus, palomino is simply a partially expressed color allele and not a set of characteristics that make up a "breed".
Because registration as a palomino with a color breed registry is based primarily on coat color, horses from many breeds or combination of breeds may qualify. Some breeds that have palomino representatives are the American Saddlebred, Tennessee Walking Horse, Morgan and Quarter Horse. The color is fairly rare in the Thoroughbred, but does in fact occur and is recognized by The Jockey Club. Some breeds, such as the Haflinger and Arabian, may appear to be palomino, but are genetically chestnuts with flaxen manes and tails, as neither breed carries the cream dilution gene. However, in spite of their lack of cream DNA, some palomino color registries have registered such horses if their coat color falls within the acceptable range of shades.
While the color standard used by palomino organizations usually describes the ideal body color as that of a "newly minted gold coin" (sometimes mistakenly claimed to be a penny), a wider a body color range is often accepted, ranging from a cream-white color to a deep, dark, chocolate color ("chocolate palomino") that may actually be silver dapple or liver chestnut with a flaxen mane and tail.
Requirements for registration
In the United States, there are two primary color breed registries for Palomino-colored horses: the Palomino Horse Association (PHA), and the Palomino Horse Breeders of America (PHBA).
The Palomino Horse Association (PHA) registers palomino horses of any breed and type "on color and conformation". The shade of color considered ideal by the PHA is the color of a gold coin, but shades of palomino from light to dark gold are accepted. The mane and tail are required to be white, silver, or ivory, but up to 15% dark or reddish-brown hair is accepted. In the interest of breeding palomino horses, the PHA also registers full double-dilute blue-eyed cremellos, erroneously called "cremello palominos" by the PHA. Horses that are not recorded by any other registry of unknown pedigree are accepted if their color meets the PHA definition of "palomino".
The Palomino Horse Breeders of America (PHBA) has stricter requirements. To be accepted by the PHBA, in addition to color, a horse must have the general structure appropriate to the breeds of light riding type recognized by the PHBA. The adult height of the PHBA horse should be , and the horse must not show draft horse or pony characteristics. An individual that does not meet the height requirements may still be accepted if it is registered in one of the breed registries recognized by the PHBA. The PHBA usually requires horses or both parents of the horse to be registered by or eligible for registration with certain recognized breed registries, including those for the American Quarter Horse, Paint, Appaloosa, Saddlebred, Morgan, Holsteiner, Arabian, assorted part-Arabian registries, Pinto (horse division only), Thoroughbred, and assorted gaited horse breeds. Horses with PHBA-registered parents are also eligible even if they are not recorded with any other breed registry. In some situations, mares and geldings may be registered without pedigree on account of their conformation and color only, but stallions must always have pedigrees that are "verified in fact".
The ideal PHBA body color is the shade of "a United States gold coin". The mane and tail must be naturally white, and may not have more than 15% black, brown or off-colored hairs. Brown or dark primitive markings are not accepted. PHBA also does not accept horses that are gray or show color characteristics of Paints, pintos, Appaloosas or cremellos or perlinos. The skin must be dark, other than pink skin on the face connected to a white marking. The PHBA will not accept a horse for regular registration if it has all three characteristics of a double-dilute cream: light (or pink) skin over the body; white or cream-colored hair over the body; and eyes of a bluish cast. White markings on the face and legs may not exceed certain limits. Leg white may not be higher than the level of the elbow or the stifle, white on the face may not extend past the throatlatch. Spotting and characteristics of the Leopard complex and the various pinto patterns are not accepted, and body spots of less than a 4-inch diameter may be allowed. Horses with non-dark skin on the body, white or creamy coat and pink skin around the eyes are not accepted. Spots of pink skin visible in the muzzle or around the eyes, under the tail and between the hind legs are not accepted. An exception is made for horses registered with the American Saddlebred Horse Association, which may have skin of any color. Accepted eye colors are black, brown, blue and hazel. However, horses with blue or partially blue eyes are accepted only if their registration certificate from a recognized breed association mentions the eye color; they are also accepted on horses of unknown pedigree if they are gelded or spayed.
| Biology and health sciences | Horses | Animals |
276054 | https://en.wikipedia.org/wiki/Hummingbird%20hawk-moth | Hummingbird hawk-moth | The hummingbird hawk-moth (Macroglossum stellatarum) is a species of hawk moth found across temperate regions of Eurasia. The species is named for its similarity to hummingbirds, as they feed on the nectar of tube-shaped flowers using their long proboscis while hovering in the air; this resemblance is an example of convergent evolution.
The hummingbird hawk-moth was first described by Carl Linnaeus in his 1758 10th edition of Systema Naturae. As of 2018, its entire genome and mitogenome have been sequenced.
Distribution
The hummingbird hawk-moth is distributed throughout the northern Old World from Portugal to Japan, but it breeds mainly in warmer climates (southern Europe, North Africa, and points east). Three generations are produced in a year in Spain. There is evidence that the population in the British Isles is actively expanding its range, as numbers have been consistently increasing. In addition, it is believed that this population is becoming resident instead of migratory, as warmer temperatures due to climate change are allowing individuals to overwinter.
It is a strong flier, dispersing widely in the summer. However it rarely survives the winter in northern latitudes (e.g. north of the Alps in Europe, north of the Caucasus in Russia).
Moths in the genus Hemaris, also of the family Sphingidae, are known as "hummingbird moths" in the US, and "bee moths" in Europe. This sometimes causes confusion between this species and the North American genus.
Life cycle
Two or more broods are produced each year. The adult may be encountered at any time of the year, especially in the south of the range, where there may be three or four broods. It overwinters as an adult in a crevice among rocks, trees, and buildings. On very warm days it may emerge to feed in mid-winter. Unlike other moths, they have no sexual dimorphism in the size of their antennal lobes.
Ova
The glossy pale green ova (eggs) are spherical with a diameter. They are said to look like the flower buds of the host plant Galium, and that is where the female lays them. They hatch 6 to 8 days after laying. Up to 200 eggs may be laid by one female, each on a separate plant.
Larvae
Newly hatched larvae are clear yellow, and in the second instar assume their green coloration. The larva is green with two grey stripes bordered in cream along the sides and with a horn at the rear end typical of sphingids. The horn is purplish red, changing to blue with an orange tip in the last instar. They feed fully exposed on the top of the host plant and rest in among a tangle of stems. Although dependent on warmth and sun, the larval stage can be as rapid as 20 days.
Pupae
The pupae are pale brownish with a prominent, keeled proboscis, and two sharp spines at the end of the cremaster. They are enclosed in loose silken cocoons among the host plant debris or on the ground among leaf litter.
Adults
The forewings are brown, with black wavy lines across them, while the hindwings are orange with a black edge. The abdomen is quite broad, with a fan-tail of setae at the end. The wingspan is .
In the southern parts of its range, the hummingbird hawk-moth is highly active even when temperatures are high, and thoracic temperatures above have been measured. This is among the highest recorded for hawk-moths, and near the limit for insect muscle activity.
Behaviour
Its long proboscis () and its hovering behavior, accompanied by an audible humming noise, make it look remarkably like a hummingbird while feeding on flowers. Like hummingbirds, it feeds on flowers which have tube-shaped corollae. It should not be confused with the moths called hummingbird moths in North America, genus Hemaris, members of the same family and with similar appearance and behavior.
The resemblance to hummingbirds is an example of convergent evolution. It flies during the day, especially in bright sunshine, but also at dusk, dawn, and even in the rain, which is unusual for even diurnal hawkmoths. M. stellatarum engages in free hovering flight, which allows more maneuverability and control than fixed-wing flight, despite high energetic cost. Like many large insects, it relies upon Johnston's organs for body positioning information.
Vision
The hummingbird hawkmoth's visual abilities have been studied extensively, and they have demonstrated a relatively good ability to learn colours. They have a trichromatic visual system, and are most sensitive to wavelength in the range of 349-521 nm. They have been shown to discriminate a wavelength difference as small as 1–2 nm between sources. This discrimination is even more precise than Apis mellifera, or the western honey bee. Among other flower visitors, their visual system is similar to Papilio xuthus, or the Asian swallowtail butterfly, and Deilephila elpenor, the nocturnal elephant hawkmoth. Their food preference is based mainly on visual identification, while D. elpenor preference relies upon olfactory identification. Compared to D. elpenor, M. stellatarum have a much smaller number of ommatidia, but a larger optic lobe volume to provide more visual processing tissue.
Habitat and host plants
Hummingbird hawk-moths can be easily seen in gardens, parks, meadows, bushes, and woodland edge, where the preferred food plants grow (honeysuckle, red valerian and many others).
Their larvae usually feed on bedstraws or madders (Rubia) but have been recorded on other Rubiaceae and Centranthus, Stellaria, and Epilobium.
Adults are particularly fond of nectar-rich flowers with a long and narrow calyx, since they can then take advantage of their long proboscis and avoid competition from other insects. Flowers with longer tubes typically present the feeding animal a higher nectar reward. Proboscis length is thought to have been evolutionarily impacted by the length of flower feeding tubes. Examples of such plants include Centranthus, Jasminum, Buddleia, Nicotiana, Primula, Viola, Syringa, Verbena, Echium, Phlox, and Stachys. They are reported to trap-line, that is, to return to the same flower beds at about the same time each day.
In culture
Hummingbird hawk-moths have been seen as a lucky omen. In particular, a swarm of the moths was seen flying across the English Channel on D-Day, the day of the Normandy landings in the Second World War. These moths, along with other moths, are in the family Sphingidae because their larvae were thought to resemble the Egyptian Sphinx. The They Might Be Giants song "The Bee of the Bird of the Moth" is about the hummingbird moth.
| Biology and health sciences | Lepidoptera | Animals |
276057 | https://en.wikipedia.org/wiki/Sphingidae | Sphingidae | The Sphingidae are a family of moths commonly called sphinx moths, also colloquially known as hawk moths, with many of their caterpillars known as hornworms. It includes about 1,450 species. It is best represented in the tropics, but species are found in every region. They are moderate to large in size and are distinguished among moths for their agile and sustained flying ability, similar enough to that of hummingbirds as to be reliably mistaken for them. Their narrow wings and streamlined abdomens are adaptations for rapid flight. The family was named by French zoologist Pierre André Latreille in 1802.
Some hawk moths, such as the hummingbird hawk-moth or the white-lined sphinx, hover in midair while they feed on nectar from flowers, so are sometimes mistaken for hummingbirds. This hovering capability is only known to have evolved four times in nectar feeders: in hummingbirds, certain bats, hoverflies, and these sphingids (an example of convergent evolution). Sphingids have been studied for their flying ability, especially their ability to move rapidly from side to side while hovering, called "swing-hovering" or "side-slipping". This is thought to have evolved to deal with ambush predators that lie in wait in flowers.
Sphingids are some of the faster flying insects; some are capable of flying at over . They have wingspans from to over .
Description
Sphingid's antennae are generally not very feathery, even in males. They lack tympanal organs, but members of the group Choerocampini have hearing organs on their heads. They have a frenulum and retinaculum to join hindwings and forewings. The thorax, abdomen, and wings are densely covered in scales. Some sphingids have a rudimentary proboscis, but most have a very long one, which is used to feed on nectar from flowers. Most are crepuscular or nocturnal, but some species fly during the day. Both males and females are relatively long-lived (10 to 30 days). Prior to flight, most species shiver their flight muscles to warm them up, and, during flight, body temperatures may surpass .
In some species, differences in form between the sexes is quite marked. For example, in the African species Agrius convolvuli (the convolvulus or morning glory hawk-moth), males have thicker antennae and more mottled wing markings than females. Only males have both an undivided frenular hook and a retinaculum. Only males have a partial comb of hairs along with their antennae. Females attract males with pheromones. The male may douse the female with a pheromone before mating.
Behavior
Some species fly only for short periods either around dusk or dawn, while other species only appear later in the evening and others around midnight, but such species may occasionally be seen feeding on flowers during the day. A few common species in Africa, such as the Oriental beehawk (Cephonodes hylas virescens), Macroglossum hirundo, and Macroglossum trochilus, are diurnal.
A number of species are known to be migratory, all in the Sphingini and Macroglossinae, and specially in the genera Agrius, Cephonodes, Macroglossum, Hippotion and Theretra.
Flight
In studies with Manduca sexta, moths have dynamic flight sensory capabilities due to their antennae. The antennae are vibrated in a plane so that when the body of the moth rotates during controlled aerial maneuvers, the antennae are subject to the inertial Coriolis forces that are linearly proportional to the angular velocity of the body. The Coriolis forces cause deflections of the antennae, which are detected by the Johnston's organ at the base of each antenna, with strong frequency responses at the beat frequency of the antennae (around 25 Hz) and at twice the beat frequency. The relative magnitude of the two frequency responses enables the moth to distinguish rotation around the different principal axes, allowing for rapid course control during aerial maneuvers.
Life cycle
Most species are multivoltine, capable of producing several generations a year if weather conditions permit. Females lay translucent, greenish, flattened, smooth eggs, usually singly on the host plants. Egg development time varies highly, from three to 21 days.
Sphingid caterpillars are medium to large in size, with stout bodies. They have five pairs of prolegs. Usually, their bodies lack any hairs or tubercules, but most species have a "horn" at the posterior end, which may be reduced to a button, or absent, in the final instar. Many are cryptic greens and browns, and have countershading patterns to conceal them. Others are more conspicuously colored, typically with white spots on a black or yellow background along the length of the body. A pattern of diagonal slashes along the side is a common feature. When resting, the caterpillar usually holds its legs off the surface and tucks its head underneath (praying position), which, resembling the Great Sphinx of Giza, gives rise to the name "sphinx moth". Some tropical larvae are thought to mimic snakes. Larvae are quick to regurgitate their sticky, often toxic, foregut contents on attackers such as ants and parasitoids. Development rate depends on temperature, and to speed development, some northern and high-altitude species sunbathe. Larvae burrow into the soil to pupate, where they remain for two to three weeks before they emerge as adults.
In some sphingids, the pupa has a free proboscis, rather than being fused to the pupal case as is most common in the macrolepidoptera. They have a cremaster at the tip of the abdomen. Usually, they pupate off the host plant, in an underground chamber, among rocks, or in a loose cocoon. In most species, the pupa is the overwintering stage.
Food plants
Larvae
Sphingid larvae tend to be specific feeders, rather than generalists. Compared to similarly sized saturniids, sphingids eat soft young leaves of host plants with small toxic molecules, and chew and mash the food into very small bits. Some species can tolerate quite high concentrations of specific toxins. Tobacco hornworms (Manduca sexta) detoxify and rapidly excrete nicotine, as do several other related sphinx moths in the subfamilies Sphinginae and Macroglossinae, but members of the Smerinthinae that were tested are susceptible. The species that are able to tolerate the toxin do not sequester it in their tissues; 98% was excreted. However, other species, such as Hyles euphorbiae and Daphnis nerii, do sequester toxins from their hosts, but do not pass them on to the adult stage.
Adults
Most adults feed on nectar, although a few tropical species feed on eye secretions, and the death's-head hawkmoth steals honey from bees. Night-flying sphingids tend to prefer pale flowers with long corolla tubes and a sweet odor, a pollination syndrome known as "sphingophily". Some species are quite general in visitations, while others are very specific, with the plant only being successfully pollinated by a particular species of moth. Orchids frequently have such specific relations with hawk moths and very long corolla tubes. The comet orchid (Angraecum sesquipedale), a rare Malagasy flower with its nectar stored at the bottom of a tube, was described in 1822 by Louis-Marie Aubert du Petit-Thouars, and later, Charles Darwin famously predicted there must be some specialized moth to feed from it:
Alfred Russel Wallace published a sort of "wanted poster" (properly, a drawing in a book) of what this lepidopteran might look like, and, concurring with his colleague, added:
The predicted sphingid was discovered 21 years later and described as a subspecies of the one African species studied by Wallace: Xanthopan morganii praedicta, for which, the subspecific name praedicta ("the predicted one") was given. The Madagascan individuals had a pink, rather than white, breast and abdomen and a black apical line on the forewing, broader than in mainland specimens. Molecular clock models using either rate- or fossil-based calibrations imply that the Madagascan subspecies X. m. praedicta and the African subspecies X. m. morgani diverged 7.4 ± 2.8 Mya (million years ago), which overlaps the divergence of A. sesquipedale from its sister, A. sororium, namely 7.5 ± 5.2 Mya. Since both these orchids have extremely long spurs, longspurs likely existed before that and were exploited by long-tongued moths similar to Xanthopan morganii praedicta. The long geological separation of subspecies morgani and praedicta matches their morphological differences in the color of the breast and abdomen.
Relationships and species
The Sphingidae is sometimes assigned its own exclusive superfamily, Sphingoidea, but is alternatively included with the more encompassing Bombycoidea. Following Hodges (1971) two subfamilies are accepted, namely the Sphinginae and Macroglossinae. Around 1,450 species of hawk moths are classified into around 200 genera. Some of the best-known hawk moth species are:
Privet hawk moth (Sphinx ligustri)
White-lined Sphinx (Hyles lineata)
Death's-head hawk moth (Acherontia atropos)
Lime hawk moth (Mimas tiliae)
Poplar hawk moth (Laothoe populi)
Convolvulus hawk moth (Agrius convolvuli)
Catalpa sphinx (Ceratomia catalpae)
Hummingbird hawk-moth (Macroglossum stellatarum)
Elephant hawk moth (Deilephila elpenor)
Vine hawk moth (Hippotion celerio)
Spurge hawk moth (Hyles euphorbiae)
Oleander hawk moth (Daphnis nerii)
Pandora sphinx moth (Eumorpha pandorus)
Tomato worm (Manduca quinquemaculata)
Tobacco hornworm (Manduca sexta)
| Biology and health sciences | Lepidoptera | Animals |
276083 | https://en.wikipedia.org/wiki/Wolf%20spider | Wolf spider | Wolf spiders are members of the family Lycosidae (), named for their robust and agile hunting skills and excellent eyesight. They live mostly in solitude, hunt alone, and usually do not spin webs. Some are opportunistic hunters, pouncing upon prey as they find it or chasing it over short distances; others wait for passing prey in or near the mouth of a burrow. Wolf spiders resemble nursery web spiders (family Pisauridae), but wolf spiders carry their egg sacs by attaching them to their spinnerets, while the Pisauridae carry their egg sacs with their chelicerae and pedipalps. Two of the wolf spider's eight eyes are large and prominent; this distinguishes them from nursery web spiders, whose eyes are all of roughly equal size. This can also help distinguish them from the similar-looking grass spiders.
Description
The many genera of wolf spiders range in body size (legs not included) from less than . They have eight eyes arranged in three rows. The bottom row consists of four small eyes, the middle row has two very large eyes (which distinguishes them from the Pisauridae), and the top row has two medium-sized eyes. Unlike most other arachnids, which are generally blind or have poor vision, wolf spiders have excellent eyesight.
The tapetum lucidum is a retroreflective tissue found in eyes. This reflective tissue is only found in four secondary eyes of the wolf spider. Flashing a beam of light over the spider produces eyeshine; this eyeshine can be seen when the lighting source is roughly coaxial with the viewer or sensor. The light from the light source (e.g., a flashlight or sunlight) has been reflected from the spider's eyes directly back toward its source, producing a "glow" that is easily noticed. Wolf spiders possess the third-best eyesight of all spider groups, bettered by jumping spiders of the family Salticidae (which can distinguish colors) and the huntsman spiders of the family Sparassidae.
Wolf spiders are unique in the way that they carry their eggs. The egg sac, a round, silken globe, is attached to the spinnerets at the end of the abdomen, allowing the spider to carry her unhatched young with her. The abdomen must be held in a raised position to keep the egg case from dragging on the ground. Despite this handicap, they are still capable of hunting. Another aspect unique to wolf spiders is their method of care of young. Immediately after the spiderlings emerge from their protective silken case, they clamber up their mother's legs and crowd onto the dorsal side of her abdomen. The mother carries the spiderlings for several weeks before they are large enough to disperse and fend for themselves.
Because they depend on camouflage for protection, they do not have the flashy appearance of some other kinds of spiders. In general, their coloration is appropriate to their favorite habitat.
Hogna is the genus with the largest of the wolf spiders. Among the Hogna species in the U.S., the nearly solid dark brown H. carolinensis (Carolina wolf spider) is the largest, with a body that can be more than long. It is sometimes confused with H. helluo, which is somewhat smaller and different in coloration. The underside of H. carolinensis is solid black, but the underside of H. helluo is variegated and has reds, oranges, and yellows with shades of black.
Some members of the Lycosidae, such as H. carolinensis, make deep, tubular burrows in which they lurk much of the time. Others, such as H. helluo, seek shelter under rocks and other shelters as nature may provide. As with spiders in general, males of almost any species can sometimes be found inside homes and buildings as they wander in search for females during the autumn.
Wolf spiders play an important role in natural population control of insects and are often considered "beneficial bugs" due to their predation of pest species within farms and gardens.
Venom
Wolf spiders inject venom if continually provoked. Symptoms of their bites include swelling and mild pain. In the past, necrotic bites have been attributed to some South American and Australian species, but further investigation has indicated that those problems that did occur were probably due to bites by members of other families or did not induce those effects.
Genera
, the World Spider Catalog accepts these genera:
Acantholycosa Dahl, 1908—Asia, Europe, North America
Adelocosa Gertsch, 1973—Hawaii
Agalenocosa Mello-Leitão, 1944—South America, Oceania, Mexico, India
Aglaoctenus Tullgren, 1905—South America
Algidus New York, 1975—USA
Allocosa Banks, 1900—Oceania, North America, Africa, South America, Costa Rica, Asia, Europe
Allotrochosina Roewer, 1960—Australia, New Zealand
Alopecosa Simon, 1885—Asia, Europe, South America, Africa, North America, Oceania
Amblyothele Simon, 1910—Africa
Anomalomma Simon, 1890—Pakistan, Indonesia, Zimbabwe
Anomalosa Roewer, 1960—Australia
Anoteropsis L. Koch, 1878—New Zealand, Papua New Guinea
Arctosa C. L. Koch, 1847—Africa, Europe, Asia, South America, North America, Vanuatu
Arctosippa Roewer, 1960—Peru
Arctosomma Roewer, 1960—Ethiopia
Artoria Thorell, 1877—Oceania, Africa, Asia
Artoriellula Roewer, 1960—South Africa, Indonesia
Artoriopsis Framenau, 2007—Australia, New Zealand
Aulonia C. L. Koch, 1847—Turkey
Auloniella Roewer, 1960—Tanzania
Birabenia Mello-Leitão, 1941—Argentina, Uruguay
Bogdocosa Ponomarev & Belosludtsev, 2008—Asia
Brevilabus Strand, 1908—Ivory Coast, Senegal, Ethiopia
Bristowiella Saaristo, 1980—Comoros, Seychelles
Camptocosa Dondale, Jiménez & Nieto, 2005—United States, Mexico
Caporiaccosa Roewer, 1960—Ethiopia
Caspicosa Ponomarev, 2007—Kazakhstan, Russia
Costacosa Framenau & Leung, 2013—Australia
Crocodilosa Caporiacco, 1947—India, Myanmar, Egypt
Cynosa Caporiacco, 1933—North Africa
Dejerosa Roewer, 1960—Mozambique
Deliriosa Kovblyuk, 2009—Ukraine
Diahogna Roewer, 1960—Australia
Diapontia Keyserling, 1877—South America
Dingosa Roewer, 1955—Australia, Peru, Brazil
Dolocosa Roewer, 1960—St. Helena
Donacosa Alderweireldt & Jocqué, 1991—Spain
Dorjulopirata Buchar, 1997—Bhutan
Draposa Kronestedt, 2010—Asia
Dzhungarocosa Fomichev & Marusik, 2017—Kazakhstan
Edenticosa Roewer, 1960—Equatorial Guinea
Evippa Simon, 1882—Africa, Asia, Spain
Evippomma Roewer, 1959—Africa, Asia
Foveosa Russell-Smith, Alderweireldt & Jocqué, 2007
Geolycosa Montgomery, 1904—Africa, South America, Asia, North America, Oceania
Gladicosa Brady, 1987—North America
Gnatholycosa Mello-Leitão, 1940—Argentina
Gulocosa Marusik, Omelko & Koponen, 2015
Hesperocosa Gertsch & Wallace, 1937—United States
Hippasa Simon, 1885—Africa, Asia
Hippasella Mello-Leitão, 1944—Argentina, Peru, Bolivia
Hoggicosa Roewer, 1960—Australia
Hogna Simon, 1885—Asia, Africa, South America, North America, Caribbean, Europe, Oceania, Central America
Hognoides Roewer, 1960—Tanzania, Madagascar
Hyaenosa Caporiacco, 1940—Asia, Africa
Hygrolycosa Dahl, 1908—Asia, Greece
Karakumosa Logunov & Ponomarev, 2020—Asia
Kangarosa Framenau, 2010—Australia
Katableps Jocqué, Russell-Smith & Alderweireldt, 2011
Knoelle Framenau, 2006—Australia
Lobizon Piacentini & Grismado, 2009—Argentina
Loculla Simon, 1910—Iran, Africa
Lycosa Latreille, 1804—North America, Africa, Caribbean, Asia, Oceania, South America, Central America, Europe
Lycosella Thorell, 1890—Indonesia
Lysania Thorell, 1890—China, Malaysia, Indonesia
Mainosa Framenau, 2006—Australia
Malimbosa Roewer, 1960—West Africa
Margonia Hippa & Lehtinen, 1983—India
Megarctosa Caporiacco, 1948—Africa, Asia, Argentina, Greece
Melecosa Marusik, Omelko & Koponen, 2015
Melocosa Gertsch, 1937—North America, Brazil
Minicosa Alderweireldt & Jocqué, 2007—South Africa
Molitorosa Roewer, 1960—Brazil
Mongolicosa Marusik, Azarkina & Koponen, 2004—Mongolia, China
Mustelicosa Roewer, 1960—Ukraine, Asia
Navira Piacentini & Grismado, 2009—Argentina
Notocosa Vink, 2002—New Zealand
Nukuhiva Berland, 1935—Marquesas Is.
Oculicosa Zyuzin, 1993—Kazakhstan, Uzbekistan, Turkmenistan
Ocyale Audouin, 1826—Africa, Peru, Asia
Orinocosa Chamberlin, 1916—South America, Africa, Asia
Ovia Sankaran, Malamel & Sebastian, 2017—India, China, Taiwan
Paratrochosina Roewer, 1960—Argentina, North America, Russia
Pardosa C. L. Koch, 1847—Asia, Europe, South America, North America, Africa, Caribbean, Oceania, Central America
Pardosella Caporiacco, 1939—Ethiopia, Tanzania
Passiena Thorell, 1890—Africa, Asia
Pavocosa Roewer, 1960—Argentina, Brazil, Thailand
Phonophilus Ehrenberg, 1831—Libya
Pirata Sundevall, 1833—South America, Africa, North America, Asia, Cuba, Central America
Piratula Roewer, 1960—Asia, North America, Ukraine
Portacosa Framenau, 2017—Australia
Proevippa Purcell, 1903—Africa
Prolycosides Mello-Leitão, 1942—Argentina
Pseudevippa Simon, 1910—Namibia
Pterartoria Purcell, 1903—South Africa, Lesotho
Pyrenecosa Marusik, Azarkina & Koponen, 2004—Europe
Rabidosa Roewer, 1960—United States
Satta Lehtinen & Hippa, 1979—Papua New Guinea
Schizocosa Chamberlin, 1904—South America, Asia, Africa, North America, Vanuatu, Central America
Shapna Hippa & Lehtinen, 1983—India
Sibirocosa Marusik, Azarkina & Koponen, 2004—Russia
Sosippus Simon, 1888—North America, Central America
Syroloma Simon, 1900—Hawaii
Tapetosa Framenau, Main, Harvey & Waldock, 2009
Tasmanicosa Roewer, 1959—Australia
Tetralycosa Roewer, 1960—Australia
Tigrosa Brady, 2012—North America
Trabea Simon, 1876—Africa, Spain, Turkey
Trabeops Roewer, 1959—North America
Trebacosa Dondale & Redner, 1981—Europe, North America
Tricassa Simon, 1910—Namibia, South Africa, Madagascar
Trochosa C. L. Koch, 1847—North America, Asia, Africa, South America, Oceania, Central America, Europe, Caribbean
Trochosippa Roewer, 1960—Africa, Indonesia, Argentina
Tuberculosa Framenau & Yoo, 2006—Australia
Varacosa Chamberlin & Ivie, 1942—North America
Venator Hogg, 1900—Australia
Venatrix Roewer, 1960—Oceania, Philippines
Venonia Thorell, 1894—Asia, Oceania
Vesubia Simon, 1910—Italy, Russia, Turkmenistan
Wadicosa Zyuzin, 1985—Africa, Asia
Xerolycosa Dahl, 1908—Asia, Tanzania
Zantheres Thorell, 1887—Myanmar
Zenonina Simon, 1898—Africa
Zoica Simon, 1898—Asia, Oceania
Zyuzicosa Logunov, 2010—Asia
Evolutionary history
Wolf spiders likely originated after the K–Pg extinction event sometime in the late Paleocene, with most main subfamilies likely originating during the Eocene and Early Oligocene between 41 and 32 million years ago.
Habitats
Wolf spiders are found in a wide range of coastal and inland habitats. These include shrublands, woodland, wet coastal forests, alpine meadows, suburban gardens, and homes. Spiderlings disperse aerially; consequently, wolf spiders have wide distributions. Although some species have very specific microhabitat needs (such as stream-side gravel beds or montane herb fields), most are wanderers without permanent homes. Some build burrows which can be left open or have a trap door (depending on species). Arid-zone species construct turrets or plug their holes with leaves and pebbles during the rainy season to protect themselves from flood waters. Often, they are found in man-made locations such as sheds and other outdoor equipment.
Mating behavior
Many species of wolf spiders possess very complex courtship behaviors and secondary sexual characteristics, such as tufts of bristles on their legs or special colorations, which are most often found on the males of the species. These sexual characteristics vary by species and are most often found as modifications of the first pair of legs. First-leg modifications are often divided into elongated bristles on the legs, increased swelling of leg segments, or the full elongation of the first pair of legs compared to the other three pairs. Some mating behaviors are common between wolf spider genera, and many more that are species-specific. In the most commonly studied genus of wolf spiders, Schizocosa, researchers found that all males engage in a seismic component of their courtship display, either stridulation, or drumming their forelegs on the ground, but some also dependent on visual cues in their courtship display, as well as the seismic signaling, such as waving the front two legs in the air in front of the female, concluding that some Schizocosa species rely on multimodal courtship behaviors.
The Lycosidae comprise mainly wandering spiders, and as such, population density and male-to-female sex ratio puts selective pressures on wolf spiders when finding mates. Female wolf spiders that have already mated are more likely to eat the next male that tries to mate with them than those that have not mated yet. Males that have already mated have a higher probability of successfully mating again, but females that have already mated have a lower probability of mating again.
Relationship to humans
Though wolf spiders do bite humans, their bites are not dangerous. Wolf spider bites often result in mild redness, itching, ulcers, and if the bite wound is not cleaned it could lead to infection. However, wolf spiders usually only bite when they feel threatened or mishandled.
Wolf spiders have been found to be a vital source of natural pest control for many people's personal gardens or even homes, since the wolf spider preys on perceived pests such as crickets, ants, cockroaches, and in some cases lizards and frogs. In recent years, wolf spiders have been utilized as pest control in agriculture to reduce the amount of pesticides needed on crops. A notable example is the use of wolf spiders in cranberry bogs as a means of controlling unwanted crop destruction.
In culture
South Carolina designated the Carolina wolf spider (Hogna carolinensis) as the official state spider in 2000 due to the efforts of Skyler B. Hutto, a third-grade student at Sheridan Elementary School in Orangeburg.
At the time, South Carolina was the only U.S. state that recognized a state spider. In 2015, efforts began to name an official state spider for neighboring North Carolina.
Gallery
| Biology and health sciences | Spiders | Animals |
276092 | https://en.wikipedia.org/wiki/Chanterelle | Chanterelle | Chanterelle is the common name of several species of fungi in the genera Cantharellus, Craterellus, Gomphus, and Polyozellus. They are orange, yellow or white, meaty and funnel-shaped. On the lower surface, most species have rounded, forked folds that run almost all the way down the stipe, which tapers down from the cap. Many species emit a fruity aroma and often have a mildly peppery taste.
Chanterelles are found in Eurasia, North America, and Africa, typically growing in forested areas. They initially gained popularity as an edible mushroom in the 18th century via their inclusion in French cuisine.
Taxonomy
At one time, all yellow or golden chanterelles in western North America had been classified as Cantharellus cibarius. Using DNA analysis, they have since been shown to be a group of related species. In 1997, the Pacific golden chanterelle (C. formosus) and C. cibarius var. roseocanus were identified, followed by C. cascadensis in 2003, C. californicus in 2008, and C. enelensis in 2017. C. cibarius var. roseocanus occurs in the Pacific Northwest in Sitka spruce forests, as well as Eastern Canada in association with Pinus banksiana.
Etymology
The name chanterelle originates from the Greek kantharos meaning "tankard" or "cup", a reference to their general shape. Its German name, Pfifferling, refers to its peppery taste.
Description
The mushrooms are orange, yellow or white, meaty and funnel-shaped. On the lower surface, underneath the smooth cap, most species have rounded, forked folds that run almost all the way down the stipe, which tapers down seamlessly from the cap. Many species emit a fruity aroma, reminiscent of apricots, and often have a mildly peppery taste.
Similar species
The false chanterelle (Hygrophoropsis aurantiaca) has a similar appearance and can be confused with the chanterelle. Distinguishing factors are that false chanterelles have true gills, while chanterelles have folds. Additionally, color can help distinguish the two; the true chanterelle is uniform egg-yellow, while the false chanterelle is more orange in hue and graded, with darker center. The true chanterelle's folds are typically more wrinkled or rounded, and randomly forked. Though once thought to be hazardous, it is now known that the false chanterelle is edible but not especially tasty, and ingesting it may result in mild gastrointestinal distress. The poisonous species in the genus Omphalotus (the jack-o'-lantern mushrooms) have been misidentified as chanterelles, but can usually be distinguished by their well-developed, unforked true gills. Species of Omphalotus are not closely related to chanterelles. Other species in the closely related genera Cantharellus and Craterellus may appear similar to the golden chanterelle.
Cantharellus pallens has sometimes been defined as a species in its own right, but it is normally considered to be just a variety (C. cibarius var. pallens). Unlike "true" C. cibarius it yellows and then reddens when touched and has a weaker smell. Eyssartier and Roux classify it as a separate species but say that 90% of the chanterelles sold in French markets are this, not C. cibarius.
Similarly, the very pale C. alborufescens, which reddens easily and is found in Mediterranean areas, and northern of Iran is sometimes distinguished as a separate variety or a separate species.
Species
An incomplete listing of species that have been called chanterelles includes:
Cantharellus cascadensis
Cantharellus cibarius, which has been split into several species
Cantharellus cinnabarinus
Cantharellus enelensis
Cantharellus formosus
Cantharellus lateritius
Cantharellus minor
Cantharellus roseocanus
Cantharellus subalbidus
Craterellus cinereus
Craterellus cornucopioides
Craterellus ignicolor
Craterellus tubaeformis
Craterellus odoratus
Gomphus clavatus
Polyozellus multiplex
Distribution and habitat
Chanterelles are common in Eurasia, North America (including Central America) and Africa. In the American Pacific Northwest, they can be found from July to November. They tend to grow in clusters in mossy coniferous forests, but are also often found in mountainous birch forests and among grasses and low-growing herbs. In central Europe, the golden chanterelle is often found in beech forests among similar species and forms. In the UK, they may be found from July through December.
Uses
Nutrition
Raw chanterelle mushrooms are 90% water, 7% carbohydrates, including 4% dietary fiber, 1.5% protein, and have negligible fat. A 100 gram reference amount of raw chanterelles supplies 38 kilocalories of food energy and the B vitamins, niacin and pantothenic acid, in rich content (20% or more of the Daily Value, DV), 27% DV of iron, with moderate contents (10-1 of riboflavin, manganese, and potassium (table).
When exposed to sunlight, raw chanterelles produce a rich amount of vitamin D2 (35% DV) – also known as ergocalciferol.
Culinary
Though records of chanterelles being eaten date back to the 16th century, they first gained widespread recognition as a culinary delicacy with the spreading influence of French cuisine in the 18th century, when they began appearing in palace kitchens. For many years, they remained notable for being served at the tables of nobility. Nowadays, the usage of chanterelles in the kitchen is common throughout Europe and North America. In 1836, the Swedish mycologist Elias Fries considered the chanterelle "as one of the most important and best edible mushrooms."
Chanterelles as a group are generally described as being rich in flavor, with a distinctive taste and aroma difficult to characterize. Some species have a fruity odor, others a more woody, earthy fragrance, and still others can even be considered spicy. The golden chanterelle is perhaps the most sought-after and flavorful chanterelle, and many chefs consider it on the same short list of gourmet fungi as truffles and morels. It therefore tends to command a high price in both restaurants and specialty stores.
There are many ways to cook chanterelles. Most of the flavorful compounds in chanterelles are fat-soluble, making them good mushrooms to sauté in butter, oil or cream. They also contain smaller amounts of water- and alcohol-soluble flavorings, which lend the mushrooms well to recipes involving wine or other cooking alcohols. Many popular methods of cooking chanterelles include them in sautés, soufflés, cream sauces, and soups. They are not typically eaten raw, as their rich and complex flavor is best released when cooked.
Chanterelles are also well-suited for drying, and tend to maintain their aroma and consistency quite well. Some chefs profess that reconstituted chanterelles are actually superior in flavor to fresh ones, though they lose in texture whatever they gain in flavor by becoming more chewy after being preserved by drying. Dried chanterelles can also be crushed into flour and used in seasoning in soups or sauces. Chanterelles are also suitable for freezing, though older frozen chanterelles can often develop a slightly bitter taste after thawing.
One mushroom guide asserts, "Chanterelles are often dirty, and when washed they soak up water like a sponge...[try] dry-sauteeing...it concentrates their flavor while allowing you to wash them."
In culture
In January 2024, the California golden chanterelle became the official mushroom of that state.
| Biology and health sciences | Edible fungi | null |
276106 | https://en.wikipedia.org/wiki/Molar%20concentration | Molar concentration | Molar concentration (also called molarity, amount concentration or substance concentration) is a measure of the concentration of a chemical species, in particular, of a solute in a solution, in terms of amount of substance per unit volume of solution. In chemistry, the most commonly used unit for molarity is the number of moles per liter, having the unit symbol mol/L or mol/dm3 in SI units. A solution with a concentration of 1 mol/L is said to be 1 molar, commonly designated as 1 M or 1 M. Molarity is often depicted with square brackets around the substance of interest; for example, the molarity of the hydrogen ion is depicted as [H+].
Definition
Molar concentration or molarity is most commonly expressed in units of moles of solute per litre of solution. For use in broader applications, it is defined as amount of substance of solute per unit volume of solution, or per unit volume available to the species, represented by lowercase :
Here, is the amount of the solute in moles, is the number of constituent particles present in volume (in litres) of the solution, and is the Avogadro constant, since 2019 defined as exactly . The ratio is the number density .
In thermodynamics, the use of molar concentration is often not convenient because the volume of most solutions slightly depends on temperature due to thermal expansion. This problem is usually resolved by introducing temperature correction factors, or by using a temperature-independent measure of concentration such as molality.
The reciprocal quantity represents the dilution (volume) which can appear in Ostwald's law of dilution.
Formality or analytical concentration
If a molecule or salt dissociates in solution, the concentration refers to the original chemical formula in solution, the molar concentration is sometimes called formal concentration or formality (FA) or analytical concentration (cA). For example, if a sodium carbonate solution () has a formal concentration of c() = 1 mol/L, the molar concentrations are c() = 2 mol/L and c() = 1 mol/L because the salt dissociates into these ions.
Units
In the International System of Units (SI), the coherent unit for molar concentration is mol/m3. However, most chemical literature traditionally uses mol/dm3, which is the same as mol/L. This traditional unit is often called a molar and denoted by the letter M, for example:
1 mol/m3 = 10−3 mol/dm3 = 10−3 mol/L = 10−3 M = 1 mM = 1 mmol/L.
The SI prefix "mega" (symbol M) has the same symbol. However, the prefix is never used alone, so "M" unambiguously denotes molar.
Sub-multiples, such as "millimolar" (mM) and "nanomolar" (nM), consist of the unit preceded by an SI prefix:
Related quantities
Number concentration
The conversion to number concentration is given by
where is the Avogadro constant.
Mass concentration
The conversion to mass concentration is given by
where is the molar mass of constituent .
Mole fraction
The conversion to mole fraction is given by
where is the average molar mass of the solution, is the density of the solution.
A simpler relation can be obtained by considering the total molar concentration, namely, the sum of molar concentrations of all the components of the mixture:
Mass fraction
The conversion to mass fraction is given by
Molality
For binary mixtures, the conversion to molality is
where the solvent is substance 1, and the solute is substance 2.
For solutions with more than one solute, the conversion is
Properties
Sum of molar concentrations – normalizing relations
The sum of molar concentrations gives the total molar concentration, namely the density of the mixture divided by the molar mass of the mixture or by another name the reciprocal of the molar volume of the mixture. In an ionic solution, ionic strength is proportional to the sum of the molar concentration of salts.
Sum of products of molar concentrations and partial molar volumes
The sum of products between these quantities equals one:
Dependence on volume
The molar concentration depends on the variation of the volume of the solution due mainly to thermal expansion. On small intervals of temperature, the dependence is
where is the molar concentration at a reference temperature, is the thermal expansion coefficient of the mixture.
Examples
| Physical sciences | Concentration | Basics and measurement |
276115 | https://en.wikipedia.org/wiki/Jumping%20spider | Jumping spider | Jumping spiders are a group of spiders that constitute the family Salticidae. , this family contained over 600 described genera and over 6,000 described species, making it the largest family of spiderscomprising 13% of spider species. Jumping spiders have some of the best vision among arthropods and use it in courtship, hunting, and navigation. Although they normally move unobtrusively and fairly slowly, most species are capable of very agile jumps, notably when hunting, but sometimes in response to sudden threats or crossing long gaps. Both their book lungs and tracheal system are well-developed, and they use both systems (bimodal breathing). Jumping spiders are generally recognized by their eye pattern. All jumping spiders have four pairs of eyes, with the anterior median pair (the two front middle eyes) being particularly large.
Description
Jumping spiders are among the easiest to distinguish from similar spider families because of the shape of the cephalothorax and their eye patterns. The families closest to Salticidae in general appearance are the Corinnidae (distinguished also by prominent spines on the back four legs), the Oxyopidae (the lynx spiders, distinguished by very prominent spines on all legs), and the Thomisidae (the crab spiders, distinguished by their front four legs, which are very long and powerful). None of these families, however, have eyes that resemble those of the Salticidae. Conversely, the legs of jumping spiders are not covered with any very prominent spines. Their front four legs generally are larger than the hind four, but not as dramatically so as those of the crab spiders, nor are they held in the outstretched-arms attitude characteristic of the Thomisidae. In spite of the length of their front legs, Salticidae depend on their rear legs for jumping. The generally larger front legs are used partly to assist in grasping prey, and in some species, the front legs and pedipalps are used in species-recognition signaling.
The jumping spiders, unlike the other families, have faces that are roughly rectangular surfaces perpendicular to their direction of motion. In effect this means that their forward-looking, anterior eyes are on "flat faces", as shown in the photographs. Their eye pattern is the clearest single identifying characteristic. They have eight eyes, as illustrated. Most diagnostic are the front row of four eyes, in which the anterior median pair are more dramatically prominent than any other spider eyes apart from the posterior median eyes of the Deinopidae. There is, however, a radical functional difference between the major (anterior median) eyes of Salticidae and the major (posterior median) eyes of the Deinopidae; the large posterior eyes of Deinopidae are adapted mainly to vision in dim light, whereas the large anterior eyes of Salticidae are adapted to detailed, three-dimensional vision for purposes of estimating the range, direction, and nature of potential prey, permitting the spider to direct its attacking leaps with great precision. The anterior lateral eyes, though large, are smaller than the anterior median eyes and provide a wider forward field of vision.
The rear row of four eyes may be described as strongly bent, or as being rearranged into two rows, with two large posterior lateral eyes being the furthest back. They serve for lateral vision. The posterior median eyes also have been shifted out laterally, almost as far as the posterior lateral eyes. They are usually much smaller than the posterior lateral eyes and there is doubt about whether they are at all functional in many species.
The body length of jumping spiders generally ranges from . The largest is Hyllus giganteus, while other genera with relatively large species include Phidippus, Philaeus and Plexippus.
In addition to using their silk for safety lines while jumping, they also build silken "pup tents", where they take shelter from bad weather and sleep at night. They molt in these shelters, build and store egg cases in them, and also spend the winter in them.
Their body's sensory hairs are able to detect airborne acoustic stimuli up to 3 m away.
Vision
Jumping spiders have four pairs of eyes; three secondary pairs that are fixed and a principal pair that is movable.
The posterior median eyes are vestigial in many species, but in some primitive subfamilies, they are comparable in size with the other secondary eyes and help to detect motion. While unable to form images, the reduced pair of eyes is thought to have a role similar to that of insect ocelli by receiving light from the sky. The photoreceptors in the other secondary pairs are almost exclusively green-sensitive, but the posterior median eyes have two visual opsins different from those in all the other eyes, sensitive to blue and UV light.
The posterior lateral eyes (PLEs) are wide-angle motion detectors that sense motions from the side and behind. Combined with the other eyes, PLEs give the spider a near 360° view of the world.
The anterior lateral eyes (ALEs) have the best visual acuity of the secondary eyes. They are able to distinguish some details, as well, and without them, no "looming response" can be triggered by motion. Even with all the other pairs covered, jumping spiders in a study could still detect, stalk, and attack flies, using their ALEs only, which are also sufficiently widely spaced to provide stereoscopic vision.
The anterior median eyes have very good vision. This pair of eyes is built like a telescopic tube with a corneal lens in the front and a second lens in the back that focus images onto a four-layered retina, a narrow, boomerang-shaped strip oriented vertically. Physiological experiments have shown they may have up to four different kinds of receptor cells, with different absorption spectra, giving them the possibility of tetrachromatic color vision, with sensitivity extending into the ultraviolet (UV) range. As the eyes are too close together to allow depth perception, and the animals do not make use of motion parallax, they have instead evolved a method called image defocus. Of the four photoreceptor layers in the retina, the two closest to the surface contain a UV-sensitive opsin (visual pigment), while the two deepest contain a green-sensitive opsin. The incoming green light is only focused on the deepest layer, while the other one receives defocused or fuzzy images. By measuring the amount of defocus from the fuzzy layer, calculating the distance to the objects in front of them is possible. In addition to receptor cells, red filters also have been detected, located in front of the cells that normally register green light. All salticids, regardless of whether they have two, three, or four kinds of color receptors, seemingly are highly sensitive to UV light. Some species (such as Cosmophasis umbratica) are highly dimorphic in the UV spectrum, suggesting a role in sexual signaling. Color discrimination has been demonstrated in behavioral experiments.
The anterior median eyes have high resolution (11 min visual angle), but the field of vision is narrow, from 2 to 5°. The central region of the retina, where acuity is highest, is no more than six or seven receptor rows wide. However, the eye can scan objects off the direct axis of vision. As the lens is attached to the carapace, the eye's scanning movements are restricted to its retina through a complicated pattern of translations and rotations. This dynamic adjustment is a means of compensation for the narrowness of the static field of vision. Movement of the retina in jumping spiders is analogous to the way many vertebrates, such as primates, move their entire eyes to focus images of interest onto their fovea centralis. In jumping spiders with a translucent carapace, such movements within the jumping spider's eyes are visible from outside when the attention of the spider is directed to various targets.
Behavior
Jumping
Many other arthropods are known to jump, including grasshoppers, fleas, leafhoppers, and sand fleas. Jumping spiders are different from these animals because they are able to make accurate, targeted jumps. Jumps are used for navigation, to escape danger, and to catch prey. When jumping, they use mainly their third or fourth pair of legs, or both pairs, depending on species. Jumping spiders' well-developed internal hydraulic system extends their limbs by altering the pressure of their body fluid (hemolymph) within them. This enables the spiders to jump without having large muscular legs like a grasshopper. The maximum horizontal jump distance varies greatly between species, with some capable of jumping two or three body lengths, while the jump of an individual Colonus puerperus was measured at 38 times the body length. The accuracy of their jumps is mediated by their well-developed visual system and the ability to quickly process visual information to tailor each jump. When a jumping spider moves from place to place, and especially just before it jumps, it tethers a filament of silk (or 'dragline') to whatever it is standing on. This dragline provides a mechanical aid to jumping, including braking and stabilization and if the jump should fail, the spider climbs back up the dragline.
Hunting
The hunting behaviour of the Salticidae is confusingly varied compared to that of most spiders in other families. Salticids hunt diurnally as a rule, which is consistent with their highly developed visual system. When it detects potential prey, a jumping spider typically begins orienting itself by swiveling its cephalothorax to bring the anterior median eyes to bear. It then moves its abdomen into line with its cephalothorax. After that, it might spend some time inspecting the object of its attention and determining whether a camouflaged or doubtful item of prey is promising, before it starts to stalk slowly forward. When close enough, the spider pauses to attach a dragline, then springs onto the prey.
Many variations on the theme and many surprising aspects exist. For one, salticids do not necessarily follow a straight path in approaching prey. They may follow a circuitous course, sometimes even a course that takes the hunter through regions from which the prey is not visible. Such complex adaptive behaviour is hard to reconcile with an organism that has such a tiny brain, but some jumping spiders, in particular some species of Portia, can negotiate long detours from one bush down to the ground, then up the stem of another bush to capture a prey item on a particular leaf. Such behaviour still is the subject of research.
Some salticid species are continually on the move, stopping periodically to look around for prey, which they then stalk immediately. Others spend more time scanning their surroundings from one position, actively stalking any prey they detect. Members of the genus Phaeacius take that strategy to extremes; they sit on a tree trunk, facing downwards and rarely do any stalking, but simply lunge down on any prey items that pass close before them.
Some Salticidae specialise in particular classes of prey, such as ants. Most spiders, including most salticids, avoid worker ants, but several species not only eat them as a primary item in their diets, but also employ specialised attack techniques; Anasaitis canosa, for example, circles around to the front of the ant and grabs it over the back of its head. Such myrmecophagous species, however, do not necessarily refuse other prey items, and routinely catch flies and similar prey in the usual salticid fashion, without the special precautions they apply in hunting dangerous prey such as ants. Ants offer the advantages of being plentiful prey items for which little competition from other predators occurs, but catching less hazardous prey when it presents itself remains profitable.
Some of the most surprising hunting behaviours occur among the araneophagous Salticidae, and vary greatly in method. Many of the spider-hunting species quite commonly attack other spiders, whether fellow salticids or not, in the same way as any other prey, but some kinds resort to web invasion; nonspecialists such as Phidippus audax sometimes attack prey ensnared in webs, basically in acts of kleptoparasitism; sometimes they leap onto and eat the web occupant itself, or simply walk over the web for that purpose.
Salticidae in the genera Brettus, Cyrba, Gelotia, and Portia display more advanced web-invasion behavior. They slowly advance onto the web and vibrate the silk with their pedipalps and legs. In this respect, their behaviour resembles that of the Mimetidae, probably the most specialised of the araneophagous spider families. If the web occupant approaches in the manner appropriate to dealing with ensnared prey, the predator attacks.
The foregoing examples present the Salticidae as textbook examples of active hunters; they would hardly seem likely to build webs other than those used in reproductive activities, and in fact, most species really do not build webs to catch prey. However, exceptions occur, though even those that do build capture webs generally also go hunting like other salticids. Some Portia species, for example, spin capture webs that are functional, though not as impressive as some orb webs of the Araneidae; Portia webs are of an unusual funnel shape and apparently adapted to the capture of other spiders. Spartaeus species, though, largely capture moths in their webs. In their review of the ethology of the Salticidae, Richman and Jackson speculate on whether such web building is a relic of the evolution of this family from web-building ancestors.
In hunting, the Salticidae also use their silk as a tether to enable them to reach prey that otherwise would be inaccessible. For example, by advancing towards the prey to less than the jumping distance, then retreating and leaping in an arc at the end of the tether line, many species can leap onto prey on vertical or even on inverted surfaces, which of course would not be possible without such a tether.
Having made contact with the prey, hunting Salticidae administer a bite to inject rapid-acting venom that gives the victim little time to react. In this respect, they resemble the Mimetidae and Thomisidae, families that ambush prey that often are larger than the predator, and they do so without securing the victim with silk; they accordingly must immobilise it immediately and their venom is adapted accordingly.
Sleeping
Among several organisms, scientists have discovered that octopuses and cuttlefish experience REM sleep. Although REM sleep has been proved to be a phase of sleep in various organisms, there was a lack of evidence that supported the presence of REM sleep in insects and arachnids.(Mason, 2022) However, in 2022, a group of researchers published an article on supporting the presence of REM sleep in what was thought to just be a normal sleep cycle in jumping spiders. The group of researchers published an article about the ways in which baby jumping spiders displayed various indicators of REM sleep, similar to those displayed by humans under this phase of deep sleep. More specifically, the jumping spiders displayed a significant amount of action occurring in their retinal tubes and uncoordinated twitches/leg curls under this state. Given that the jumping spiders were observed from 7PM to 7AM, the researchers realized that these actions were not present while they stretched or readjusted their silk webs outside of this time frame. Ultimately, solidifying the indicators of REM sleep in these small organisms.(Rößler et al., 2022)
Diet
Although jumping spiders are generally carnivorous, many species have been known to include nectar in their diets, and one species, Bagheera kiplingi, feeds primarily on plant matter. None are known to feed on seeds or fruit. Extrafloral nectaries on plants, such as Chamaecrista fasciculata (partridge pea), provide jumping spiders with nectar; the plant benefits accordingly when the spiders prey on whatever pests they find.
The female of the Southeast Asian species Toxeus magnus feeds its offspring with a milky, nutritious fluid for the first 40 days of their lives.
Reproduction
Courtship and mating behavior
Jumping spiders conduct complex, visual courtship displays using movements and physical bodily attributes. A form of sexual dimorphism, the males possess plumose hairs, colored or iridescent hairs (particularly pronounced in the peacock spiders), front leg fringes, structures on other legs, and other, often bizarre, modifications. These characteristics are used in a courtship "dance" in which the colored or iridescent parts of the body are displayed. In addition to displaying colors, jumping spiders perform complex sliding, vibrational, or zigzag movements to attract females. Many males have auditory signals, as well. These amplified sounds presented to the females resemble buzzes or drum rolls. Species vary significantly in visual and vibratory components of courtship. The ability to sense UV light (see Vision section) is used by at least one species, Cosmophasis umbratica, in courtship behavior, though it is reasonable to assume that many other species exhibit this characteristic. Cosmophasis umbratica males have markings that are only visible in UV and the females use the markings for mate choice.
If receptive to the male, the female assumes a passive, crouching position. In some species, the female may vibrate her palps or abdomen. The male then extends his front legs towards the female to touch her. If the female remains receptive, the male climbs on her back and inseminates her with his palps.
Consequences of sexual dimorphism
Maintaining colorful ornamentation may seem strictly beneficial to sexual selection, yet costs to maintain such distinguishing characteristics occur. While colorful or UV-reflecting individuals may attract more female spiders, it can also increase the risk of predation.
Taxonomy
The monophyly of the family Salticidae is well established through both phylogenetic and morphological analyses. The sister group to Salticidae is the family Philodromidae. Synapomorphies of the two families include loss of cylindrical gland spigots and loss of tapeta in the indirect eyes.
A 2015 revision of the Salticidae family divided it into seven subfamilies:
Onomastinae Maddison, 2015 – 1 extant genus
Asemoneinae Maddison, 2015 – 4 extant genera (Hindumanes, originally placed here, was moved to Lyssomaninae)
Lyssomaninae Blackwall, 1877 – 4 extant genera (including Hindumanes)
Spartaeinae Wanless, 1984 – 29 extant genera in 3 tribes
Eupoinae Maddison, 2015 – 3 extant genera
Hisponinae Simon, 1901 – 6 extant genera
Salticinae Blackwall, 1841 – about 540 extant genera in 27 tribes
The relationships between these subfamilies is still up for debate. Below are the results of a 2017 phylogenomic study that attempted to resolve this question. The subfamily Eupoinae was unevaluated and its exact position is unclear.
Habitat
Jumping spiders live in a variety of habitats. Tropical forests harbor the most species, but they are also found in temperate forests, scrubland, deserts, intertidal zones, and mountainous regions. Euophrys omnisuperstes is the species reported to have been collected at the highest elevation, on the slopes of Mount Everest.
Models for mimicry
Some small insects are thought to have evolved an appearance or behavioural traits that resemble those of jumping spiders and this is suspected to prevent their predation, specifically from jumping spiders. Some examples appear to be provided by patterns on the wings of some tephritid flies, the nymph of a fulgorid and possibly some moths.
Fossils
Very few jumping spider fossils have been found. Of those known, all are from Cenozoic era amber. The oldest fossils are from Baltic amber dating to the Eocene epoch, specifically, 54 to 42 million years ago. Other fossil jumping spiders have been preserved within Chiapas amber and Dominican amber.
| Biology and health sciences | Spiders | Animals |
276294 | https://en.wikipedia.org/wiki/Goldschmidt%20classification | Goldschmidt classification | The Goldschmidt classification,
developed by Victor Goldschmidt (1888–1947), is a geochemical classification which groups the chemical elements within the Earth according to their preferred host phases into lithophile (rock-loving), siderophile (iron-loving), chalcophile (sulfide ore-loving or chalcogen-loving), and atmophile (gas-loving) or volatile (the element, or a compound in which it occurs, is liquid or gaseous at ambient surface conditions).
Some elements have affinities to more than one phase. The main affinity is given in the table below and a discussion of each group follows that table.
Lithophile elements
Lithophile elements () are those that remain on or close to the surface because they combine readily with oxygen, forming compounds that did not sink into the Earth's core. The lithophile elements include Al, B, Ba, Be, Br, Ca, Cl, Cr, Cs, F, I, Hf, K, Li, Mg, Na, Nb, O, P, Rb, Sc, Si, Sr, Ta, Th, Ti, U, V, Y, Zr, W and the lanthanides or rare earth elements (REE).
Lithophile elements mainly consist of the highly reactive metals of the s- and f-blocks. They also include a small number of reactive nonmetals, and the more reactive metals of the d-block such as titanium, zirconium and vanadium.
Most lithophile elements form very stable ions with an electron configuration of a noble gas (sometimes with additional f-electrons). The few that do not, such as silicon, phosphorus and boron, form strong covalent bonds with oxygen, often involving pi bonding. Their strong affinity for oxygen causes lithophile elements to associate very strongly with silica, forming relatively low-density minerals that thus rose towards the crust during planetary differentiation. The more soluble minerals formed by the alkali metals tend to concentrate in seawater or arid regions where they can crystallise. The less soluble lithophile elements are concentrated on ancient continental shields where soluble minerals have been weathered.
Because of their strong affinity for oxygen, most lithophile elements are enriched in the Earth's crust relative to their abundance in the solar system. The most reactive s- and f-block metals, which form either saline or metallic hydrides, are known to be extraordinarily enriched on Earth as a whole relative to their solar abundances. This is because during the earliest stages of the Earth's formation, the abundance of stable forms of each element was determined by how readily it forms volatile hydrides; these volatiles then could "escape" the proto-Earth, leaving behind those elements unreactive with hydrogen. Under these conditions, the s- and f-block metals were strongly enriched during the formation of the Earth. The most enriched elements are rubidium, strontium and barium, which between them account for over 50 percent by mass of all elements heavier than iron in the Earth's crust.
The nonmetallic lithophilesphosphorus and the halogensexist on Earth as ionic salts with s-block metals in pegmatites and seawater. With the exception of fluorine, whose hydride forms hydrogen bonds and is therefore of relatively low volatility, these elements have had their concentrations on Earth significantly reduced through escape of volatile hydrides during the Earth's formation. Although they are present in the Earth's crust in concentrations quite close to their solar abundances, phosphorus and the heavier halogens are probably significantly depleted on Earth as a whole relative to their solar abundances.
Several transition metals, including chromium, molybdenum, iron and manganese, show both lithophile and siderophile characteristics and can be found in both these two layers. Although these metals form strong bonds with oxygen and are never found in the Earth's crust in the free state, metallic forms of these elements are thought very likely to exist in the core of the earth as relics from when the atmosphere did not contain oxygen. Like the "pure" siderophiles, these elements (except iron) are considerably depleted in the crust relative to their solar abundances.
Owing to their strong affinity for oxygen, lithophile metals, although they form the great bulk of the metallic elements in Earth's crust, were never available as free metals before the development of electrolysis. With this development, many lithophile metals are of considerable value as structural metals (magnesium, aluminium, titanium, vanadium) or as reducing agents (sodium, magnesium, calcium).
The non-metals phosphorus and the halogens were also not known to early chemists, though production of these elements is less difficult than of metallic lithophiles since electrolysis is required only with fluorine. Elemental chlorine is particularly important as an oxidizing agentusually being made by electrolysis of sodium chloride.
Siderophile elements
Siderophile elements () are the transition metals which tend to sink towards the core during planetary differentiation, because they dissolve readily in iron either as solid solutions or in the molten state. Some sources include elements which are not transition metals in their list of siderophiles, such as germanium. Other sources may also differ in their list based on the temperature being discussedniobium, vanadium, chromium, and manganese may be considered siderophiles or not, depending on the assumed temperature and pressure. Also confusing the issue is that some elements, such as the aforementioned manganese, as well as molybdenum, form strong bonds with oxygen, but in the free state (as they existed on the early Earth when free oxygen did not exist) can mix so easily with iron that they do not concentrate in the siliceous crust, as do true lithophile elements. Iron, meanwhile, is simply everywhere.
The siderophile elements include the highly siderophilic ruthenium, rhodium, palladium, rhenium, osmium, iridium, platinum, and gold, the moderately siderophilic cobalt and nickel, in addition to the "disputed" elements mentioned earliersome sources even include tungsten and silver.
Most siderophile elements have practically no affinity for oxygen: indeed oxides of gold are thermodynamically unstable. They form stronger bonds with carbon or sulfur, but even these are not strong enough to separate out with the chalcophile elements. Thus, siderophile elements are bound with iron through metallic bonding in the Earth's core, where pressures may be high enough to keep the iron solid. Manganese, iron, and molybdenum do form strong bonds with oxygen, but in the free state (as on the early Earth) can mix so easily with iron that they do not concentrate in the siliceous crust, as do true lithophile elements. However, ores of manganese are found in much the same sites as are those of aluminium and titanium, owing to manganese's great reactivity towards oxygen.
Because they are so concentrated in the dense core, siderophile elements are known for their rarity in the Earth's crust. Most of them have always been known as precious metals because of this. Iridium is the rarest transition metal occurring within the Earth's crust, with an abundance by mass of less than one part per billion. Mineable deposits of precious metals usually form as a result of the erosion of ultramafic rocks, but are not highly concentrated even compared to their crustal abundances, which are typically several orders of magnitude below their solar abundances. However, because they are concentrated in the Earth's mantle and Earth's core, siderophile elements are believed to be present in the Earth as a whole (including the core) in something approaching their solar abundances.
Chalcophile elements
The chalcophile elements () include Ag, As, Bi, Cd, Cu, Ga, Ge, Hg, In, Pb, S, Sb, Se, Sn, Te, Tl and Zn.
Chalcophile elements are those that remain on or close to the surface because they combine readily with sulfur and some other chalcogens other than oxygen, forming compounds which did not sink along with iron towards the Earth's core. Chalcophile elements are those metals and heavier nonmetals that have a low affinity for oxygen and prefer to bond with sulfur as highly insoluble sulfides.
Because these sulfides are much denser than the silicate minerals formed by lithophile elements, chalcophile elements separated below the lithophiles at the time of the first crystallization of the Earth's crust. This has led to their depletion in the Earth's crust relative to their solar abundances, though because the minerals they form are nonmetallic, this depletion has not reached the levels found with siderophile elements.
However, because they formed volatile hydrides in the accreting protosolar nebula when the controlling redox reaction was the oxidation or reduction of hydrogen, the less metallic chalcophile elements are strongly depleted on Earth as a whole relative to cosmic abundances. This is most especially true of the chalcogens selenium and tellurium (which formed volatile hydrogen selenide and hydrogen telluride, respectively), which for this reason are among the rarest elements found in the Earth's crust (to illustrate, tellurium is only about as abundant as platinum).
The most metallic chalcophile elements (of the copper, zinc and boron groups) may mix to some degree with iron in the Earth's core. They are not likely to be depleted on Earth as a whole relative to their solar abundances since they do not form volatile hydrides. Zinc and gallium are somewhat "lithophile" in nature because they often occur in silicate or related minerals and form quite strong bonds with oxygen. Gallium, notably, is sourced mainly from bauxite, an aluminum hydroxide ore in which gallium ions substitute for chemically similar aluminum.
Although no chalcophile element is of high abundance in the Earth's crust, chalcophile elements constitute the bulk of commercially important metals. This is because, whereas lithophile elements require energy-intensive electrolysis for extraction, chalcophiles can be easily extracted by reduction, and chalcophiles' geochemical concentrationwhich in extreme cases can exceed 100,000 times their average crustal abundance. These greatest enrichments occur in high plateaus like the Tibetan Plateau and the Bolivian Altiplano where large quantities of chalcophile elements have been uplifted through plate tectonics. A side-effect of this in modern times is that the rarest chalcophiles (like mercury) are so completely exploited that their value as minerals has almost completely disappeared.
Atmophile elements
The atmophile elements () are H, C, N and the noble gases.
Atmophile elements (also called "volatile elements") are defined as those that remain mostly on or above Earth's surface because they are, or occur in, liquids and/or gases at temperatures and pressures found on the surface. The noble gases do not form stable compounds and occur as monatomic gases, while nitrogen, although highly reactive as the free atom, bonds so strongly into diatomic molecular nitrogen that all oxides of nitrogen are thermodynamically unstable with respect to nitrogen and oxygen. Consequently, with the development of free oxygen in Earth's atmosphere, ammonia was oxidised to molecular nitrogen which has come to form four-fifths of the Earth's atmosphere. Carbon is also classed as an atmophile because it forms very strong multiple bonds with oxygen in carbon monoxide (slowly oxidised in the atmosphere) and carbon dioxide. The latter is the fourth-largest constituent of the Earth's atmosphere, while carbon monoxide occurs naturally from various sources (volcanoes, combustion) and has a residence time in the atmosphere of a few months.
Hydrogen, which occurs in water, is also classed as an atmophile. Water is classified as a volatile, because most of it is liquid or gas, even though it can exist as a solid compound at Earth's surface. Water can also be incorporated into other minerals as water of crystallization (as in gypsum) or through ionic and hydrogen bonding (as in talc), giving hydrogen some lithophile character.
Because all atmophile elements are either gases or form volatile hydrides, atmophile elements are strongly depleted on Earth as a whole relative to their solar abundances owing to losses from the atmosphere during the formation of the Earth. The heavier noble gases (krypton, xenon) are the rarest stable elements on Earth. (In fact they, along with neon, were all first isolated and described by William Ramsay and Morris Travers and assistants, who gave them names with Ancient Greek derivations of 'hidden', 'stranger', and 'new', respectively.)
Argon is the exception among the noble gases: it is the third-most abundant component of Earth's present-day atmosphere after nitrogen and oxygen, comprising . Argon-40 is a stable daughter of radioactive potassium-40, and argon is heavy enough to be gravitationally captured by the post-accretion Earth, so while the proto-Earth's primordial argon was mostly driven off, this radiogenic argon has accumulated over geologic time. This makes Earth's argon abundance substantially different from cosmic abundance ratios for argon, being enormously enriched in , while predominates cosmically.
Trace and synthetic elements
Synthetic elements are excluded from the classification, as they do not occur naturally.
Trace radioactive elements (namely Tc, Pm, Po, At, Rn, Fr, Ra, Ac, Pa, Np, Pu) are also treated as synthetic. Although these do occur in nature, their occurrence is dependent on their long-lived parents Th and U, and they are not very mobile. For instance, polonium's chemistry would predict it to be a chalcophile, but it tends to occur instead as a lithophile along with its parent uranium. Even radon, a gas at standard conditions, does not usually have time to travel very far from the original uranium source before decaying. When needed, these elements are typically produced synthetically in nuclear reactors instead of extraction from ores.
| Physical sciences | Geochemistry | Earth science |
276409 | https://en.wikipedia.org/wiki/Raspberry | Raspberry | The raspberry is the edible fruit of several plant species in the genus Rubus of the rose family, most of which are in the subgenus Idaeobatus. The name also applies to these plants themselves. Raspberries are perennial with woody stems.
World production of raspberries in 2022 was 947,852 tonnes, led by Russia with 22% of the total. Raspberries are cultivated across northern Europe and North America and are consumed in various ways, including as whole fruit and in preserves, cakes, ice cream, and liqueurs. Raspberries are a rich source of vitamin C, manganese, and dietary fiber.
Description
A raspberry is an aggregate fruit, developing from the numerous distinct carpels of a single flower. What distinguishes the raspberry from its blackberry relatives is whether or not the torus (receptacle or stem) "picks with" (i.e., stays with) the fruit. When picking a blackberry fruit, the torus stays with the fruit. With a raspberry, the torus remains on the plant, leaving a hollow core in the raspberry fruit.
Raspberries are grown for the fresh fruit market and for commercial processing into individually quick frozen (IQF) fruit, purée, juice, or dried fruit used in a variety of grocery products such as raspberry pie. Raspberries need ample sun and water for optimal development. Raspberries thrive in well-drained soil with a pH between 6 and 7 with ample organic matter to assist in retaining water. While moisture is essential, wet and heavy soils or excess irrigation can bring on Phytophthora root rot, which is one of the most serious pest problems faced by the red raspberry. As a cultivated plant in moist, temperate regions, it is easy to grow and tends to spread unless pruned. Escaped raspberries frequently appear as garden weeds, spread by seeds found in bird droppings.
An individual raspberry weighs and is made up of around 100 drupelets, each of which consists of a juicy pulp and a single central seed. A raspberry bush can yield several hundred berries a year.
Etymology
Raspberry derives its name from raspise, "a sweet rose-colored wine" (mid-15th century), from the Anglo-Latin vinum raspeys, or from raspoie, meaning "thicket", of Germanic origin. The name may have been influenced by its appearance as having a rough surface, related to the Old English rasp or "rough berry".
Species
Examples of raspberry species in Rubus subgenus Idaeobatus include:
Rubus crataegifolius (Asian raspberry)
Rubus gunnianus (Tasmanian alpine raspberry)
Rubus idaeus (red raspberry or European red raspberry)
Rubus leucodermis (whitebark raspberry or western raspberry, blue raspberry, black raspberry)
Rubus occidentalis (black raspberry)
Rubus parvifolius (Australian native raspberry)
Rubus phoenicolasius (wine raspberry or wineberry)
Rubus rosifolius (Mauritius raspberry)
Rubus strigosus (American red raspberry) (syn. R. idaeus var. strigosus)
Rubus ellipticus (yellow Himalayan raspberry)
Rubus fraxinifolius (mountain raspberry, in southeast Asia and the Pacific Ocean)
Several species of Rubus, also called raspberries, are classified in other subgenera, including:
Rubus deliciosus (boulder raspberry, subgenus Anoplobatus)
Rubus odoratus (flowering raspberry, subgenus Anoplobatus)
Rubus nivalis (snow raspberry, subgenus Chamaebatus)
Rubus arcticus (Arctic raspberry, subgenus Cyclactis)
Rubus sieboldii (Molucca raspberry, subgenus Malachobatus)
Cultivation
Various kinds of raspberries can be cultivated from hardiness zones 3 to 9. Raspberries are traditionally planted in the winter as dormant canes, although planting of tender, plug plants produced by tissue culture has become much more common. A specialized production system called "long cane production" involves growing canes for a year in a northern climate such as Scotland or Oregon or Washington, where the chilling requirement for proper bud break is attained, or attained earlier than the ultimate place of planting. These canes are then dug, roots and all, to be replanted in warmer climates such as Spain, where they quickly flower and produce a very early season crop. Plants are typically planted 2–6 per m in fertile, well drained soil; raspberries are usually planted in raised beds/ridges, if there is any question about root rot problems.
All cultivars of raspberries have perennial roots, but many do not have perennial shoots. In fact, most raspberries have shoots that are biennial (meaning shoots grow in the first growing season and fruits grow on those shoots during the second growing season). The flowers can be a major nectar source for honeybees and other pollinators.
Raspberries are vigorous and can be locally invasive. They propagate using basal shoots (also known as suckers), extended underground shoots that develop roots and individual plants. They can sucker new canes some distance from the main plant. For this reason, raspberries spread well, and can take over gardens if left unchecked. Raspberries are often propagated using cuttings, and will root readily in moist soil conditions.
The fruit is harvested when it comes off the receptacle easily and has turned a deep color (red, black, purple, or golden yellow, depending on the species and cultivar). This is when the fruits are ripest and sweetest.
High tunnel bramble production offers the opportunity to bridge gaps in availability during late fall and late spring. Furthermore, high tunnels allow less hardy floricane-fruiting raspberries to overwinter in climates where they would not otherwise survive. In the tunnel, plants are established at close spacing usually prior to tunnel construction.
Cultivars
Major cultivars
Raspberries are an important commercial fruit crop, widely grown in all temperate regions of the world. Many of the most important modern commercial red raspberry cultivars derive from hybrids between R. idaeus and R. strigosus. Some botanists consider the Eurasian and American red raspberries to belong to a single, circumboreal species, Rubus idaeus, with the European plants then classified as either R. idaeus subsp. idaeus or R. idaeus var. idaeus, and the native North American red raspberries classified as either R. idaeus subsp. strigosus, or R. idaeus var. strigosus. Recent breeding has resulted in cultivars that are thornless and more strongly upright, not needing staking.
The black raspberry, Rubus occidentalis, is also cultivated, providing both fresh and frozen fruit, as well as jams, preserves, and other products, all with that species' distinctive flavor.
Purple raspberries have been produced by horticultural hybridization of red and black raspberries, and have also been found in the wild in a few places (for example, in Vermont) where the American red and the black raspberries both grow naturally. Commercial production of purple-fruited raspberries is rare.
Blue raspberry is a local name used in Prince Edward County, Ontario, Canada, for the cultivar 'Columbian', a hybrid (purple raspberry) of R. strigosus and R. occidentalis. Blue raspberry can also refer to the whitebark raspberry, R. leucodermis.
Both the red and the black raspberry species have albino-like pale-yellow natural or horticultural variants, resulting from presence of recessive genes that impede production of anthocyanin pigments. Fruits from such plants are called golden raspberries or yellow raspberries; despite their similar appearance, they retain the distinctive flavor of their respective species (red or black). Most pale-fruited raspberries commercially sold in the eastern United States are derivatives of red raspberries. Yellow-fruited variants of the black raspberry are sometimes grown in home gardens.
Red raspberries have also been crossed with various species in other subgenera of the genus Rubus, resulting in a number of hybrids, the first of which was the loganberry. Later notable hybrids include the olallieberry, boysenberry, marionberry, and tayberry; all are multi-generational hybrids. Hybridization between the familiar cultivated red raspberries and a few Asiatic species of Rubus has also been achieved.
Selected cultivars
Numerous raspberry cultivars have been selected.
Two types of raspberry are available for commercial and domestic cultivation; the summer-bearing type produces an abundance of fruit on second-year canes (floricanes) within a relatively short period in midsummer, and double or "everbearing" plants, which also bear some fruit on first-year canes (primocanes) in the late summer and fall, as well as the summer crop on second-year canes. Those marked (AGM) have gained the Royal Horticultural Society's Award of Garden Merit.
Red, early Summer fruiting
Boyne
Cascade Dawn
Fertödi Venus
Glen Clova
Glen Moy (AGM)
Killarney
Latham
Malahat
Malling Exploit
Malling Jewel (AGM)
Prelude
Rubin Bulgarski
Titan
Willamette
Red, Mid-summer Fruiting
Cuthbert
Glen Ample (AGM)
Lloyd George
Meeker
Newburgh
Ripley
Skeena
Cowichan
Chemainus
Saanich
Red, Late Summer Fruiting
Cascade Delight
Coho
Fertödi Rubina
Glen Magna (AGM)
Leo (AGM)
Malling Admiral (AGM)
Octavia
Schoenemann
Tulameen (AGM)
Red primocane, Autumn fruiting
Amity
Augusta
Autumn Bliss (AGM)
Joan J. (Thornless)
Caroline
Fertödi Kétszertermö
Heritage
Imara
Joan J
Josephine
Kwanza
Kweli
Mapema
Polka (AGM)
Rafiki
Ripley
Summit
Zeva Herbsternte
Yellow primocane, Autumn fruiting
Anne
Fallgold
Fertödi Aranyfürt
Goldenwest
Golden Queen
Honey Queen
Jambo
Kiwi Gold
Purple (hybrids between black and red raspberries)
Brandywine
Glencoe
Royalty
Black
Black Hawk
Bristol
Cumberland
Jewel
Logan
Morrison
Munger
Ohio Everbearer
Scepter
Dwarf cultivars
= 'Nr7'
Diseases and pests
Raspberries are sometimes eaten by the larvae of some Lepidoptera species (butterflies and moths). More serious are the raspberry beetle (in Europe) and the raspberry fruitworm (in North America), whose larvae can damage raspberries.
Botrytis cinerea, or gray mold, is a common fungal infection of raspberries and other soft fruit under wet conditions. It is seen as a gray mold growing on the raspberries, and particularly affects fruit which are bruised, as the bruises provide an easy entrance point for the spores.
Raspberry plants should not be planted where potatoes, tomatoes, peppers, eggplants, or bulbs have previously been grown, without prior fumigation of the soil. These crops are hosts for the disease Verticillium wilt, a fungus that can stay in the soil for many years and can infest the raspberry crop.
Animals
Raspberries, among other plants with high sugar content like peaches, are prime targets for the Japanese beetle, which relies heavily on these sources as its main food resource. The voracious feeding habits of Japanese beetles not only pose a direct threat to raspberry plants but also increase the risk of transmitting various plant diseases. This dual impact can significantly undermine agricultural productivity, making it crucial for raspberry growers to implement effective pest management strategies to mitigate the damage caused by Japanese beetle infestations.
Production
In 2022, world production of raspberries was 947,852 tonnes, led by Russia with 22% of the total (table). Other major producers were Mexico, Serbia, Poland, and the United States.
Nutrition
Raw raspberries are 86% water, 12% carbohydrates, and have about 1% each of protein and fat (table). In a reference amount of , raspberries supply 53 kilocalories and 6.5 grams of dietary fiber. Raspberries are a rich source of vitamin C (29% of the Daily Value, DV), manganese (29% DV), and dietary fiber (26% DV), but otherwise have low content of micronutrients (table). Raspberries are a low-glycemic index food, with total sugar content of only 4% and no starch.
The aggregate fruit structure contributes to raspberry's nutritional value, as it increases the proportion of dietary fiber, which is among the highest known in whole foods up to 6% fiber per total weight.
Phytochemicals
Raspberries contain phytochemicals, such as anthocyanin pigments, ellagic acid, ellagitannins, quercetin, gallic acid, cyanidins, pelargonidins, catechins, kaempferol and salicylic acid. Yellow raspberries and others with pale-colored fruits are lower in anthocyanin content. Both yellow and red raspberries contain carotenoids, mostly lutein esters, but these are masked by anthocyanins in red raspberries.
Raspberry compounds are under preliminary research for their potential to affect human health.
Leaves
Raspberry leaves can be used fresh or dried in herbal teas, providing an astringent flavor. In herbal and traditional medicine, raspberry leaves are used for some remedies, although there is no scientifically valid evidence to support their medicinal use.
| Biology and health sciences | Rosales | null |
276410 | https://en.wikipedia.org/wiki/Module%20%28mathematics%29 | Module (mathematics) | In mathematics, a module is a generalization of the notion of vector space in which the field of scalars is replaced by a (not necessarily commutative) ring. The concept of a module also generalizes the notion of an abelian group, since the abelian groups are exactly the modules over the ring of integers.
Like a vector space, a module is an additive abelian group, and scalar multiplication is distributive over the operations of addition between elements of the ring or module and is compatible with the ring multiplication.
Modules are very closely related to the representation theory of groups. They are also one of the central notions of commutative algebra and homological algebra, and are used widely in algebraic geometry and algebraic topology.
Introduction and definition
Motivation
In a vector space, the set of scalars is a field and acts on the vectors by scalar multiplication, subject to certain axioms such as the distributive law. In a module, the scalars need only be a ring, so the module concept represents a significant generalization. In commutative algebra, both ideals and quotient rings are modules, so that many arguments about ideals or quotient rings can be combined into a single argument about modules. In non-commutative algebra, the distinction between left ideals, ideals, and modules becomes more pronounced, though some ring-theoretic conditions can be expressed either about left ideals or left modules.
Much of the theory of modules consists of extending as many of the desirable properties of vector spaces as possible to the realm of modules over a "well-behaved" ring, such as a principal ideal domain. However, modules can be quite a bit more complicated than vector spaces; for instance, not all modules have a basis, and, even for those that do (free modules), the number of elements in a basis need not be the same for all bases (that is to say that they may not have a unique rank) if the underlying ring does not satisfy the invariant basis number condition, unlike vector spaces, which always have a (possibly infinite) basis whose cardinality is then unique. (These last two assertions require the axiom of choice in general, but not in the case of finite-dimensional vector spaces, or certain well-behaved infinite-dimensional vector spaces such as Lp spaces.)
Formal definition
Suppose that R is a ring, and 1 is its multiplicative identity.
A left R-module M consists of an abelian group and an operation such that for all r, s in R and x, y in M, we have
,
,
,
The operation · is called scalar multiplication. Often the symbol · is omitted, but in this article we use it and reserve juxtaposition for multiplication in R. One may write RM to emphasize that M is a left R-module. A right R-module MR is defined similarly in terms of an operation .
Authors who do not require rings to be unital omit condition 4 in the definition above; they would call the structures defined above "unital left R-modules". In this article, consistent with the glossary of ring theory, all rings and modules are assumed to be unital.
An (R,S)-bimodule is an abelian group together with both a left scalar multiplication · by elements of R and a right scalar multiplication ∗ by elements of S, making it simultaneously a left R-module and a right S-module, satisfying the additional condition for all r in R, x in M, and s in S.
If R is commutative, then left R-modules are the same as right R-modules and are simply called R-modules.
Examples
If K is a field, then K-modules are called K-vector spaces (vector spaces over K).
If K is a field, and K[x] a univariate polynomial ring, then a K[x]-module M is a K-module with an additional action of x on M by a group homomorphism that commutes with the action of K on M. In other words, a K[x]-module is a K-vector space M combined with a linear map from M to M. Applying the structure theorem for finitely generated modules over a principal ideal domain to this example shows the existence of the rational and Jordan canonical forms.
The concept of a Z-module agrees with the notion of an abelian group. That is, every abelian group is a module over the ring of integers Z in a unique way. For , let (n summands), , and . Such a module need not have a basis—groups containing torsion elements do not. (For example, in the group of integers modulo 3, one cannot find even one element that satisfies the definition of a linearly independent set, since when an integer such as 3 or 6 multiplies an element, the result is 0. However, if a finite field is considered as a module over the same finite field taken as a ring, it is a vector space and does have a basis.)
The decimal fractions (including negative ones) form a module over the integers. Only singletons are linearly independent sets, but there is no singleton that can serve as a basis, so the module has no basis and no rank, in the usual sense of linear algebra. However this module has a torsion-free rank equal to 1.
If R is any ring and n a natural number, then the cartesian product Rn is both a left and right R-module over R if we use the component-wise operations. Hence when , R is an R-module, where the scalar multiplication is just ring multiplication. The case yields the trivial R-module {0} consisting only of its identity element. Modules of this type are called free and if R has invariant basis number (e.g. any commutative ring or field) the number n is then the rank of the free module.
If Mn(R) is the ring of matrices over a ring R, M is an Mn(R)-module, and ei is the matrix with 1 in the -entry (and zeros elsewhere), then eiM is an R-module, since . So M breaks up as the direct sum of R-modules, . Conversely, given an R-module M0, then M0⊕n is an Mn(R)-module. In fact, the category of R-modules and the category of Mn(R)-modules are equivalent. The special case is that the module M is just R as a module over itself, then Rn is an Mn(R)-module.
If S is a nonempty set, M is a left R-module, and MS is the collection of all functions , then with addition and scalar multiplication in MS defined pointwise by and , MS is a left R-module. The right R-module case is analogous. In particular, if R is commutative then the collection of R-module homomorphisms (see below) is an R-module (and in fact a submodule of NM).
If X is a smooth manifold, then the smooth functions from X to the real numbers form a ring C∞(X). The set of all smooth vector fields defined on X forms a module over C∞(X), and so do the tensor fields and the differential forms on X. More generally, the sections of any vector bundle form a projective module over C∞(X), and by Swan's theorem, every projective module is isomorphic to the module of sections of some vector bundle; the category of C∞(X)-modules and the category of vector bundles over X are equivalent.
If R is any ring and I is any left ideal in R, then I is a left R-module, and analogously right ideals in R are right R-modules.
If R is a ring, we can define the opposite ring Rop, which has the same underlying set and the same addition operation, but the opposite multiplication: if in R, then in Rop. Any left R-module M can then be seen to be a right module over Rop, and any right module over R can be considered a left module over Rop.
Modules over a Lie algebra are (associative algebra) modules over its universal enveloping algebra.
If R and S are rings with a ring homomorphism , then every S-module M is an R-module by defining . In particular, S itself is such an R-module.
Submodules and homomorphisms
Suppose M is a left R-module and N is a subgroup of M. Then N is a submodule (or more explicitly an R-submodule) if for any n in N and any r in R, the product (or for a right R-module) is in N.
If X is any subset of an R-module M, then the submodule spanned by X is defined to be where N runs over the submodules of M that contain X, or explicitly , which is important in the definition of tensor products of modules.
The set of submodules of a given module M, together with the two binary operations + (the module spanned by the union of the arguments) and ∩, forms a lattice that satisfies the modular law:
Given submodules U, N1, N2 of M such that , then the following two submodules are equal: .
If M and N are left R-modules, then a map is a homomorphism of R-modules if for any m, n in M and r, s in R,
.
This, like any homomorphism of mathematical objects, is just a mapping that preserves the structure of the objects. Another name for a homomorphism of R-modules is an R-linear map.
A bijective module homomorphism is called a module isomorphism, and the two modules M and N are called isomorphic. Two isomorphic modules are identical for all practical purposes, differing solely in the notation for their elements.
The kernel of a module homomorphism is the submodule of M consisting of all elements that are sent to zero by f, and the image of f is the submodule of N consisting of values f(m) for all elements m of M. The isomorphism theorems familiar from groups and vector spaces are also valid for R-modules.
Given a ring R, the set of all left R-modules together with their module homomorphisms forms an abelian category, denoted by R-Mod (see category of modules).
Types of modules
Finitely generated An R-module M is finitely generated if there exist finitely many elements x1, ..., xn in M such that every element of M is a linear combination of those elements with coefficients from the ring R.
Cyclic A module is called a cyclic module if it is generated by one element.
Free A free R-module is a module that has a basis, or equivalently, one that is isomorphic to a direct sum of copies of the ring R. These are the modules that behave very much like vector spaces.
Projective Projective modules are direct summands of free modules and share many of their desirable properties.
Injective Injective modules are defined dually to projective modules.
Flat A module is called flat if taking the tensor product of it with any exact sequence of R-modules preserves exactness.
Torsionless A module is called torsionless if it embeds into its algebraic dual.
Simple A simple module S is a module that is not {0} and whose only submodules are {0} and S. Simple modules are sometimes called irreducible.
Semisimple A semisimple module is a direct sum (finite or not) of simple modules. Historically these modules are also called completely reducible.
Indecomposable An indecomposable module is a non-zero module that cannot be written as a direct sum of two non-zero submodules. Every simple module is indecomposable, but there are indecomposable modules that are not simple (e.g. uniform modules).
Faithful A faithful module M is one where the action of each in R on M is nontrivial (i.e. for some x in M). Equivalently, the annihilator of M is the zero ideal.
Torsion-free A torsion-free module is a module over a ring such that 0 is the only element annihilated by a regular element (non zero-divisor) of the ring, equivalently implies or .
Noetherian A Noetherian module is a module that satisfies the ascending chain condition on submodules, that is, every increasing chain of submodules becomes stationary after finitely many steps. Equivalently, every submodule is finitely generated.
Artinian An Artinian module is a module that satisfies the descending chain condition on submodules, that is, every decreasing chain of submodules becomes stationary after finitely many steps.
Graded A graded module is a module with a decomposition as a direct sum over a graded ring such that for all x and y.
Uniform A uniform module is a module in which all pairs of nonzero submodules have nonzero intersection.
Further notions
Relation to representation theory
A representation of a group G over a field k is a module over the group ring k[G].
If M is a left R-module, then the action of an element r in R is defined to be the map that sends each x to rx (or xr in the case of a right module), and is necessarily a group endomorphism of the abelian group . The set of all group endomorphisms of M is denoted EndZ(M) and forms a ring under addition and composition, and sending a ring element r of R to its action actually defines a ring homomorphism from R to EndZ(M).
Such a ring homomorphism is called a representation of R over the abelian group M; an alternative and equivalent way of defining left R-modules is to say that a left R-module is an abelian group M together with a representation of R over it. Such a representation may also be called a ring action of R on M.
A representation is called faithful if and only if the map is injective. In terms of modules, this means that if r is an element of R such that for all x in M, then . Every abelian group is a faithful module over the integers or over some ring of integers modulo n, Z/nZ.
Generalizations
A ring R corresponds to a preadditive category R with a single object. With this understanding, a left R-module is just a covariant additive functor from R to the category Ab of abelian groups, and right R-modules are contravariant additive functors. This suggests that, if C is any preadditive category, a covariant additive functor from C to Ab should be considered a generalized left module over C. These functors form a functor category C-Mod, which is the natural generalization of the module category R-Mod.
Modules over commutative rings can be generalized in a different direction: take a ringed space (X, OX) and consider the sheaves of OX-modules (see sheaf of modules). These form a category OX-Mod, and play an important role in modern algebraic geometry. If X has only a single point, then this is a module category in the old sense over the commutative ring OX(X).
One can also consider modules over a semiring. Modules over rings are abelian groups, but modules over semirings are only commutative monoids. Most applications of modules are still possible. In particular, for any semiring S, the matrices over S form a semiring over which the tuples of elements from S are a module (in this generalized sense only). This allows a further generalization of the concept of vector space incorporating the semirings from theoretical computer science.
Over near-rings, one can consider near-ring modules, a nonabelian generalization of modules.
| Mathematics | Algebra | null |
276504 | https://en.wikipedia.org/wiki/Engineered%20wood | Engineered wood | Engineered wood, also called mass timber, composite wood, man-made wood, or manufactured board, includes a range of derivative wood products which are manufactured by binding or fixing the strands, particles, fibres, or veneers or boards of wood, together with adhesives, or other methods of fixation to form composite material. The panels vary in size but can range upwards of and in the case of cross-laminated timber (CLT) can be of any thickness from a few inches to or more. These products are engineered to precise design specifications, which are tested to meet national or international standards and provide uniformity and predictability in their structural performance. Engineered wood products are used in a variety of applications, from home construction to commercial buildings to industrial products. The products can be used for joists and beams that replace steel in many building projects. The term mass timber describes a group of building materials that can replace concrete assemblies.
Typically, engineered wood products are made from the same hardwoods and softwoods used to manufacture lumber. Sawmill scraps and other wood waste can be used for engineered wood composed of wood particles or fibers, but whole logs are usually used for veneers, such as plywood, medium-density fibreboard (MDF), or particle board. Some engineered wood products, like oriented strand board (OSB), can use trees from the poplar family, a common but non-structural species.
Alternatively, it is also possible to manufacture similar engineered bamboo from bamboo; and similar engineered cellulosic products from other lignin-containing materials such as rye straw, wheat straw, rice straw, hemp stalks, kenaf stalks, or sugar cane residue, in which case they contain no actual wood but rather vegetable fibers.
Flat-pack furniture is typically made out of man-made wood due to its low manufacturing costs and its low weight.
Types of products
There are a wide variety of engineered wood products for both structural and non-structural applications. This list is not comprehensive, and is intended to help categorize and distinguish between different types of engineered wood.
Wood-based panels
Wood-based panels are made from fibres, flakes, particles, veneers, chips, sawdust, slabs, wood powder, strands, or other wood derivate through a binding process with adhesives. Wood structural panels are a collection of flat panel products, used extensively in building construction for sheathing, decking, cabinetry and millwork, and furniture. Examples include plywood and oriented strand board (OSB). Non-structural wood-based panels are flat-panel products, used in non-structural construction applications and furniture. Non-structural panels are usually covered with paint, wood veneer, or resin paper in their final form. Examples include fiberboard and particle board.
Plywood
Plywood, a wood structural panel, is sometimes called the original engineered wood product. Plywood is manufactured from sheets of cross-laminated veneer and bonded under heat and pressure with durable, moisture-resistant adhesives. By alternating the grain direction of the veneers from layer to layer, or "cross-orienting", panel strength and stiffness in both directions are maximized. Other structural wood panels include oriented strand boards and structural composite panels.
Oriented strand board
Oriented strand board (OSB) is a wood structural panel manufactured from rectangular-shaped strands of wood that are oriented lengthwise and then arranged in layers, laid up into mats, and bonded together with moisture-resistant, heat-cured adhesives. The individual layers can be cross-oriented to provide strength and stiffness to the panel. Similar to plywood, most OSB panels are delivered with more strength in one direction. The wood strands in the outermost layer on each side of the board are normally aligned into the strongest direction of the board. Arrows on the product will often identify the strongest direction of the board (the height, or longest dimension, in most cases). Produced in huge, continuous mats, OSB is a solid panel product of consistent quality with no laps, gaps, or voids. OSB is delivered in various dimensions, strengths, and levels of water resistance.
OSB and plywood are often used interchangeably in building construction.
Fibreboard
Medium-density fibreboard (MDF) and high-density fibreboard (hardboard or HDF) are made by breaking down hardwood or softwood residuals into wood fibers, combining them with wax and a resin binder, and forming panels by applying high temperature and pressure. MDF is used in non-structural applications.
Particle board
Particle board is manufactured from wood chips, sawmill shavings, or even sawdust, and a synthetic resin or another suitable binder, which is pressed and extruded. Research published in 2017 showed that durable particle board can be produced from agricultural waste products, such as rice husk or guinea corn husk. Particleboard is cheaper, denser, and more uniform than conventional wood and plywood and is substituted for them when the cost is more important than strength and appearance. A major disadvantage of particleboard is that it is very prone to expansion and discoloration due to moisture, particularly when it is not covered with paint or another sealer. Particle board is used in non-structural applications.
Structural composite lumber
Structural composite lumber (SCL) is a class of materials made with layers of veneers, strands, or flakes bonded with adhesives. Unlike wood structural panels, structural composite lumber products generally have all grain fibers oriented in the same direction. The SCL family of engineered wood products are commonly used in the same structural applications as conventional sawn lumber and timber, including rafters, headers, beams, joists, rim boards, studs, and columns. SCL products have higher dimensional stability and increased strength compared to conventional lumber products.
Laminated veneer
Laminated veneer lumber (LVL) is produced by bonding thin wood veneers together in a large billet, similar to plywood. The grain of all veneers in the LVL billet is parallel to the long direction (unlike plywood). The resulting product features enhanced mechanical properties and dimensional stability that offer a broader range in product width, depth, and length than conventional lumber.
Parallel-strand
Parallel-strand lumber (PSL) consists of long veneer strands laid in parallel formation and bonded together with an adhesive to form the finished structural section. The length-to-thickness ratio of strands in PSL is about 300. A strong, consistent material, it has a high load-carrying ability and is resistant to seasoning stresses so it is well suited for use as beams and columns for post and beam construction, and for beams, headers, and lintels for light framing construction.
Laminated strand
Laminated strand lumber (LSL) and oriented strand lumber (OSL) are manufactured from flaked wood strands that have a high length-to-thickness ratio. Combined with an adhesive, the strands are oriented and formed into a large mat or billet and pressed. LSL and OSL offer good fastener-holding strength and mechanical-connector performance and are commonly used in a variety of applications, such as beams, headers, studs, rim boards, and millwork components. LSL is manufactured from relatively short strands—typically about long—compared to the strands used in PSL. The length-to-thickness ratio of strands is about 150 for LSL and 75 for OSL.
I-joists
I-joists are ""-shaped structural members designed for use in floor and roof construction. An I-joist consists of top and bottom flanges of various widths united with webs of various depths. The flanges resist common bending stresses, and the web provides shear performance. I-joists are designed to carry heavy loads over long distances while using less lumber than a dimensional solid wood joist of a size necessary to do the same task. As of 2004, approximately 81% of all wood light framed floors were framed using I-joists.
Mass timber
Mass timber, also known as engineered timber, is a class of large structural wood components for building construction. Mass timber components are made of lumber or veneers bonded with adhesives or mechanical fasteners. Certain types of mass timber, such as nail-laminated timber and glue-laminated timber, have existed for over a hundred years. Mass timber enjoyed increasing popularity from 2012 onward, due to growing concern around the sustainability of building materials, and interest in prefabrication, off site construction, and modularization, for which mass timber is well suited. The various types of mass timber share the advantage of faster construction times as the components are manufactured off-site, and pre-finished to exact dimensions for simple on-site fastening. Mass timber has been shown to have structural properties competitive with steel and concrete, opening the possibility to build large, tall buildings out of wood. Extensive testing has demonstrated the natural fire resistance properties of mass timber primarily due the creation of a char layer around a column or beam which prevents fire from reaching the inner layers of wood. In recognition of the proven structural and fire performance of mass timber, the International Building Code, a model code that forms the basis of many North American building codes, adopted new provisions in the 2021 code cycle that permit mass timber to be used in high-rise construction up to 18 stories.
Cross-laminated timber
Cross-laminated timber (CLT) is a versatile multi-layered panel made of lumber. Each layer of boards is placed perpendicular to adjacent layers for increased rigidity and strength. It is relatively new and gaining popularity within the construction industry as it can be used for long spans and all assemblies, e.g. floors, walls, or roofs.
Glued laminated timber
Glued laminated timber (glulam) is composed of several layers of dimensional timber glued together with moisture-resistant adhesives, creating a large, strong, structural member that can be used as vertical columns or horizontal beams. Glulam can also be produced in curved shapes, offering extensive design flexibility.
Dowel-laminated timber
Dowel laminated timber (DLT), sometimes referred to as Brettstapel, is a wood-on-wood timber. The biggest benefit of this method is that no glue or metal is needed, thus eliminating VOCs (such as formaldehyde) associated with wood adhesives used in most other engineered timbers.
Similar to CLT, DLT uses a cross laminated pattern with softwoods, but instead of wood adhesives to fix lumbers in place, holes are drilled vertically or in a 45° angle, and 15-20mm dowels made of dry hardwood or densified wood (such as thermal-compressed) are placed between the lumbers.
As the hardwood dowel absorbs moisture from the softwood to reach an equilibrium moisture content, it expands into the surrounding wood, creating a connection and 'locking' them together through friction. The dowels can be dried (such as through a kiln) prior to fitting, to maximize their expansion.
Nail-laminated timber
Nail laminated timber (NLT) is a mass timber product that consists of parallel boards fastened with nails. It can be used to create floors, roofs, walls, and elevator shafts within a building. It is one of the oldest types of mass timber, being used in warehouse construction during the Industrial Revolution. Like DLT, no chemical adhesives are used, and wood fibers are oriented in the same direction.
Engineered wood flooring
Engineered wood flooring is a type of flooring product, similar to hardwood flooring, made of layers of wood or wood-based composite laminated together. The floor boards are usually milled with a tongue-and-groove profile on the edges for consistent joinery between boards.
Lamella
The lamella is the face layer of the wood that is visible when installed. Typically, it is a sawn piece of timber. The timber can be cut in three different styles: flat-sawn, quarter-sawn, and rift-sawn.
Types of core/substrate
Wood ply construction ("sandwich core"): Uses multiple thin plies of wood adhered together. The wood grain of each ply runs perpendicular to the ply below it. Stability is attained from using thin layers of wood that have little to no reaction to climatic change. The wood is further stabilized due to equal pressure being exerted lengthwise and widthwise from the plies running perpendicular to each other.
Finger core construction: Finger core engineered wood floors are made of small pieces of milled timber that run perpendicular to the top layer (lamella) of wood. They can be 2-ply or 3-ply, depending on their intended use. If it is three-ply, the third ply is often plywood that runs parallel to the lamella. Stability is gained through the grains running perpendicular to each other, and the expansion and contraction of wood are reduced and relegated to the middle ply, stopping the floor from gapping or cupping.
Fibreboard: The core is made up of medium or high-density fibreboard. Floors with a fibreboard core are hygroscopic and must never be exposed to large amounts of water or very high humidity - the expansion caused by absorbing water combined with the density of the fibreboard, will cause it to lose its form. Fibreboard is less expensive than timber and can emit higher levels of harmful gases due to its relatively high adhesive content.
An engineered flooring construction that is popular in parts of Europe is the hardwood lamella, softwood core laid perpendicular to the lamella, and a final backing layer of the same noble wood used for the lamella. Other noble hardwoods are sometimes used for the back layer but must be compatible. This is thought by many to be the most stable of engineered floors.
Other types of modified wood
Techniques have been introduced in the field of engineered wood including transformation of natural wood in laboratories through chemical and/or physical treatments to achieve tailored mechanical, optical, thermal, and conduction properties.
Densified wood
Densified wood can be made by using a mechanical hot press to compress wood fibers, sometimes in combination with chemical modification of the wood. These processes have been shown to increase the density by a factor of three. This increase in density is expected to enhance the strength and stiffness of the wood by a proportional amount. Studies published in 2018 combined chemical processes with traditional mechanical hot press methods. These chemical processes break down lignin and hemicellulose that are found naturally in the wood. Following dissolution, the cellulose strands that remain are mechanically hot compressed. Compared to the three-fold increase in strength observed from hot pressing alone, chemically processed wood has been shown to yield an 11-fold improvement. This extra strength comes from hydrogen bonds formed between the aligned cellulose nanofibers.
The densified wood possessed mechanical strength properties on par with steel used in building construction, opening the door for applications of densified wood in situations where regular strength wood would fail. Environmentally, wood requires significantly less carbon dioxide to produce than steel.
Synthetic resin densified wood is resin-impregnated densified wood, also known as compreg. Usually phenolic resin is used as impregnation resin to impregnate and laminate plywood layers. Sometimes layers are not impregnated before lamination. It is also possible to impregnate wood chips to produce molded pressed wood components.
Delignified wood
Removing lignin from wood has several other applications, apart from providing structural advantages. Delignification alters the mechanical, thermal, optical, fluidic and ionic properties and functions of the natural wood and is an effective approach to regulating its thermal properties, as it removes the thermally conductive lignin component, while generating a large number of nanopores in the cell walls which help reduce temperature change. Delignified wood reflects most incident light and appears white in color. White wood (also known as nanowood) has high reflection haze, as well as high emissivity in the infrared wavelengths. These two characteristics generate a passive radiative cooling effect, with an average cooling power of over a 24-hour period, meaning that this wood does not "absorb" heat and therefore only emits the heat embedded in it. Moreover, white wood not only possesses a lower thermal conductivity than natural wood, and it has better thermal performance than most commercially available insulating materials. The modification of the mesoporous structure of the wood is responsible for the changes in wood performance.
White wood can also be put through a compression process, similar to the process mentioned for densified wood, which increases its mechanical performance compared to natural wood (8.7 times higher in tensile strength and 10 times higher in toughness). The thermal and structural advantages of nanowood make it an attractive material for energy-efficient building construction. However, the changes made in the wood's structural properties, like the increase in structural porosity and the partially isolated cellulose nanofibrils, damage the material's mechanical robustness. To deal with this issue, several strategies have been proposed, with one being to further densify the structure, and another to use cross-linking. Other suggestions include hybridizing natural wood with other organic particles and polymers to enhance its thermal insulation performance.
Moldable wood
Using similar chemical modification techniques to chemically densified wood, wood can be made extremely moldable using a combination of delignification and water shock treatment. This is an emerging technology and is not yet used in industrial processes. However, initial tests show promising advantages in improved mechanical properties, with the molded wood exhibiting strength comparable to some metal alloys.
Transparent wood composites
Transparent wood composites are new materials, as of 2020 are made at the laboratory scale, that combines transparency and stiffness via a chemical process that replaces light-absorbing compounds, such as lignin, with a transparent polymer.
Environmental benefits
New construction is in high demand due to growing worldwide population. However, the main materials used in new construction are currently steel and concrete. The manufacturing of these materials creates comparatively high emissions of carbon dioxide () into the atmosphere. Engineered wood has the potential to reduce carbon emissions if it replaces steel and/or concrete in the construction of buildings.
In 2014, steel and cement production accounted for about 1320 megatonnnes (Mt) and 1740 Mt respectively, which made up about 9% of global emissions that year. In a study that did not take the carbon sequestration potential of engineered wood into account, it was found that roughly 50 Mt e (carbon dioxide equivalent) could be eliminated by 2050 with the full uptake of a hybrid construction system utilizing engineered wood and steel. When considering the added effects that carbon sequestration can have over the lifetime of the material, the emissions reductions of engineered wood is even more substantial, as laminated wood that is not incinerated at the end of its lifecycle absorbs around 582 kg of /m3, while reinforced concrete emits 458 kg /m3 and steel 12.087 kg /m3.
There is not a strong consensus for measuring the carbon sequestration potential of wood. In life-cycle assessment, sequestered carbon is sometimes called biogenic carbon. ISO 21930, a standard that governs life cycle assessment, requires the biogenic carbon from a wood product can only be included as a negative input (i.e. carbon sequestration) when the wood product originated in a sustainably managed forest. This generally means that wood needs to be FSC or SFI-certified to qualify as carbon sequestering.
Advantages
Engineered wood products are used in a variety of ways, often in applications similar to solid wood products:
Mass timber (MT) is lightweight allowing the material to be easily handled, manufactured, and transported. This contributes to it being cost effective and easy to use on site.
MT offers greater strength and stiffness (based on its strength to weight ratio), increased dimensional stability, and uniformity in structures than solid wood.
When compared to steel/concrete, MT built buildings use up to 15% less energy because of the reduced energy needed to create these wood products.
MT buildings on average save 20-25% in time when compared to conventional steel/concrete buildings and 4.2% on capital cost.
MT products sequester carbon and store it within themselves over their lifespan. Using this instead of concrete and steel in buildings will reduce the embodied emissions in buildings.
Using MT has an estimated savings of around 20% in embodied carbon when compared to steel or concrete. This is because MT is a lot lighter when compared to these two materials, so it is less intensive for the machinery to transport both to site and once delivered.
MT products can provide high levels of airtightness and low coefficients of thermal conductivity meaning that the air inside cannot escape, and heat is not lost easily.
MT built buildings perform very well in seismic events because they are roughly half the mass and half the stiffness when compared to reinforced concrete buildings. Reduced stiffness allows MT buildings to be ductile and resist lateral distortion without compromising the structural integrity of the building.
MT is fire resistant to an extent. Although it is considered a combustible material, MT burns slowly and in a predictable manner. When burned, a charred layer forms on the outside that protects the inner layers of the material. However, if the charred layer comes off, the inner layers will be exposed which can compromise the integrity of the material.
Advantages by product type:
CLT: Offers high dimensional stability, high strength and stiffness and is easy to manufacture.
Glulam: Offers high strength and stiffness, is structurally efficient, and can be manufactured into complex shapes.
NLT: Does not require any specialized equipment to manufacture, is cost effective, and easy to handle.
DLT: Offers high dimensional stability, is easy and safe to manufacture, and no metal fasteners or adhesive is required.
SCL: Is able to withstand greater loads compared to solid timber and is not prone to shrinking, splitting or warping.
Engineered wood products may be preferred over solid wood in some applications due to certain comparative advantages:
Because engineered wood is man-made, it can be designed to meet application-specific performance requirements. Required shapes and dimension do not drive source tree requirements (length or width of the tree)
Engineered wood products are versatile and available in a wide variety of thicknesses, sizes, grades, and exposure durability classifications, making the products ideal for use in unlimited construction, industrial, and home project application.
Engineered wood products are designed and manufactured to maximize the natural strength and stiffness characteristics of wood. The products are very stable and some offer greater structural strength than typical wood building materials.
Glued laminated timber (glulam) has greater strength and stiffness than comparable dimensional lumber and, pound for pound, is stronger than steel.
Engineered wood panels are easy to work with using ordinary tools and basic skills. They can be cut, drilled, routed, jointed, glued, and fastened. Plywood can be bent to form curved surfaces without loss of strength. Large panel sizes speeds up construction by reducing the number of pieces that need to be handled and installed.
Engineered wood products are a more efficient use of wood as they can be made from wood that has defects, underutilized species or smaller pieces of wood which also enables the use of smaller trees
Wooden trusses are competitive in many roof and floor applications, and their high strength-to-weight ratios permit long spans offering flexibility in floor layouts.
Sustainable design advocates recommend using engineered wood, which can be produced from relatively small trees, rather than large pieces of solid dimensional lumber, which requires cutting a large tree.
Disadvantages
Like solid wood, when exposed to high moisture conditions or termites, biodeteriorations and/or fungi decay will occur which reduces the structural integrity and durability of the wood product; essentially the wood will start to rot.
Potential widespread deforestation without a sustainable forestry management plan.
MT buildings are susceptible to wind driven oscillation because of the relative flexibility of the MT material which may cause discomfort to people in the building.
Disadvantages by product type:
CLT and Glulam: More costly than solid wood.
NLT: Labor-intensive to make with potential for human error.
DLT: Limited panel sizing and thickness.
SCL: Limited panel sizing and thickness; more suitable for low rise buildings.
When compared to solid wood the following disadvantages are prevalent:
They require more primary energy for their manufacture than solid lumber.
The adhesives used in some products may cause harmful emissions. A concern with some resins is the release of formaldehyde in the finished product, often seen with urea-formaldehyde bonded products.
Properties
Plywood and OSB typically have a density of . For example, plywood sheathing or OSB sheathing typically has a surface density of . Many other engineered woods have densities much higher than OSB.
Adhesives
The types of adhesives used in engineered wood include:
Urea-formaldehyde resins (UF): most common, cheapest, and not waterproof.
Phenol formaldehyde resins (PF): yellow/brown, and commonly used for exterior exposure products.
Melamine-formaldehyde resins (MF): white, heat, and water-resistant, and often used in exposed surfaces in more costly designs.
Polymeric methylene diphenyl diisocyanate (pMDI) or polyurethane (PU) resins: expensive, generally waterproof, and does not contain formaldehyde, notoriously more difficult to release from platens and engineered wood presses.
A more inclusive term is structural composites. For example, fiber cement siding is made of cement and wood fiber, while cement board is a low-density cement panel, often with added resin, faced with fiberglass mesh.
Health concerns
While formaldehyde is an essential ingredient of cellular metabolism in mammals, studies have linked prolonged inhalation of formaldehyde gases to cancer. Engineered wood composites have been found to emit potentially harmful amounts of formaldehyde gas in two ways: unreacted free formaldehyde and the chemical decomposition of resin adhesives. When excessive amounts of formaldehyde are added to a process, the surplus will not have any additive to bond with and may seep from the wood product over time. Cheap urea-formaldehyde (UF) adhesives are largely responsible for degraded resin emissions. Moisture degrades the weak UF molecules, resulting in potentially harmful formaldehyde emissions. McLube offers release agents and platen sealers designed for those manufacturers who use reduced-formaldehyde UF and melamine-formaldehyde adhesives. Many OSB and plywood manufacturers use phenol-formaldehyde (PF) because phenol is a much more effective additive. Phenol forms a water-resistant bond with formaldehyde that will not degrade in moist environments. PF resins have not been found to pose significant health risks due to formaldehyde emissions. While PF is an excellent adhesive, the engineered wood industry has started to shift toward polyurethane binders like pMDI to achieve even greater water resistance, strength, and process efficiency. pMDIs are also used extensively in the production of rigid polyurethane foams and insulators for refrigeration. pMDIs outperform other resin adhesives, but they are notoriously difficult to release and cause buildup on tooling surfaces.
Mechanical fasteners
Some engineered wood products, such as DLT, NLT, and some brands of CLT, can be assembled without the use of adhesives using mechanical fasteners or joinery. These can range from profiled interlocking jointed boards, proprietary metal fixings, nails or timber dowels.
Building codes and standards
Throughout the years mass timber was used in buildings, codes were added to and adopted by the International Building Code (IBC) to create standards for them for the proper use and handling. For example, in 2015, CLT was incorporated into the IBC. The 2021 IBC is the latest issue of building codes, and has added three new codes regarding construction with timber material. The new three construction types go as follows, IV-A, IV-B, and IV-C, and they allow mass timber to be used in buildings up to 18, 12, and nine stories respectively.
The following technical performance standards are related to engineered wood products:
EN 300 - Oriented Strand Boards (OSB) — Definitions, classification, and specifications
EN 309 - Particleboards — Definition and classification
EN 338 - Structural timber - Strength classes
EN 386 - Glued laminated timber — performance requirements and minimum production requirements
EN 313-1 - Plywood — Classification and terminology Part 1: Classification
EN 313-2 - Plywood — Classification and terminology Part 2: Terminology
EN 314-1 - Plywood — Bonding quality — Part 1: Test methods
EN 314-2 - Plywood — Bonding quality — Part 2: Requirements
EN 315 - Plywood — Tolerances for dimensions
EN 387 - Glued laminated timber — large finger joints - performance requirements and minimum production requirements
EN 390 - Glued laminated timber — sizes - permissible deviations
EN 391 - Glued laminated timber — shear test of glue lines
EN 392 - Glued laminated timber — Shear test of glue lines
EN 408 - Timber structures — Structural timber and glued laminated timber — Determination of some physical and mechanical properties
EN 622-1 - Fibreboards — Specifications — Part 1: General requirements
EN 622-2 - Fibreboards — Specifications — Part 2: Requirements for hardboards
EN 622-3 - Fibreboards — Specifications — Part 3: Requirements for medium boards
EN 622-4 - Fibreboards — Specifications — Part 4: Requirements for soft boards
EN 622-5 - Fibreboards — Specifications — Part 5: Requirements for dry process boards (MDF)
EN 1193 - Timber structures — Structural timber and glued laminated timber - Determination of shear strength and mechanical properties perpendicular to the grain
EN 1194 - Timber structures — Glued laminated timber - Strength classes and determination of characteristic values
EN 1995-1-1 - Eurocode 5: Design of timber structures — Part 1-1: General — Common rules and rules for buildings
EN 12369-1 - Wood-based panels — Characteristic values for structural design — Part 1: OSB, particleboards, and fibreboards
EN 12369-2 - Wood-based panels — Characteristic values for structural design — Part 2: Plywood
EN 12369-3 - Wood-based panels — Characteristic values for structural design — Part 3: Solid wood panels
EN 14080 - Timber structures — Glued laminated timber — Requirements
EN 14081-1 - Timber structures - Strength graded structural timber with rectangular cross-section - Part 1: General requirements
The following product category rules can be used to create environmental product declarations for engineered wood products:
EN 15804 - Sustainability of construction works - Environmental product declarations - Core rules for the product category of construction products
EN 16485 - Round and sawn timber - Environmental Product Declarations - Product category rules for wood and wood-based products for use in construction (complementary-PCR to EN 15804)
ISO 21930 - Sustainability in buildings and civil engineering works - Core rules for environmental product declarations of construction products and services
Examples of mass timber structures
Plyscrapers
Plyscrapers are skyscrapers that are either partially made of wood or entirely made of wood. Around the world, many different plyscrapers have been built, including the Ascent MKE building, Mjostarnet in Norway, and the Stadthaus building.
The Ascent MKE building was built in 2022 in Milwaukee, Wisconsin, and is the tallest high-rise building using different mass timber components in combination with some steel and concrete. This plyscraper is 87 meters tall and has 25 stories.
The Stadthaus building is a residential building built in 2009 in Hackney, London. It has 9 stories reaching 30 meters tall. It uses CLT panels as load-bearing walls and floor 'slabs'.
The Black & White Building is an office building topped out in 2023 in Shoreditch, London. It has 6 stories reaching 17.8 meters tall. It uses CLT panels, glulam curtain walling, and LVL columns and beams.
As of 2022, over 84 mass timber buildings at least eight stories tall were in construction or completed worldwide, with numerous other projects in the planning stages. Its environmental benefits and distinctive appearance drive the growing interest in mass timber construction.
Bridges
The Mistissini Bridge built in Quebec, Canada, in 2014 is a 160-meter-long bridge that features both glulam beams and CLT panels. The bridge was designed to cross over the Uupaachikus Pass.
The Placer River Pedestrian Bridge built in Alaska, United States, in 2013. It spans long and is located in the Chugach National Forest. This bridge features glulam as it was used create the trusses.
Parking structures
The Glenwood CLT Parking Garage in Springfield, Oregon, is going to be a garage that features CLT. It will be 4 stories tall and hold 360 parking spaces. The parking garage however is under construction , and the year of completion is not yet known.
| Technology | Building materials | null |
276544 | https://en.wikipedia.org/wiki/Ground%20sloth | Ground sloth | Ground sloths are a diverse group of extinct sloths in the mammalian superorder Xenarthra. They varied widely in size with the largest, belonging to genera Lestodon, Eremotherium and Megatherium, being around the size of elephants. Ground sloths represent a paraphyletic group, as living tree sloths are thought to have evolved from ground sloth ancestors.
The early evolution of ground sloths took place during the late Paleogene and Neogene of South America, while the continent was isolated. At their earliest appearance in the fossil record, they were already distinct at the family level. Sloths dispersed into the Greater Antilles during the Oligocene, and the presence of intervening islands between the American continents in the Miocene allowed a dispersal of some species into North America. They were hardy as evidenced by their high species diversity and their presence in a wide variety of environments, extending from the far south of Patagonia (Cueva del Milodón Natural Monument) to Alaska. Sloths, and xenarthrans as a whole, represent one of the more successful South American groups during the Great American Interchange after the connection of North and South America during the late Pliocene with a number of ground sloth genera migrating northwards. One genus, Thalassocnus, even adapted for marine life along the Pacific coast of South America during the late Miocene and Pliocene epochs.
Ground sloths, which were represented by over 30 living species during the Late Pleistocene, abruptly became extinct on the American mainland as part of the end-Pleistocene extinction event around 12,000 years ago, simultaneously along with the majority of other large animals in the Americas. Their extinction has been posited to be the result of hunting by recently arrived humans and/or climate change. A number of kill sites are known where humans butchered ground sloths dating just prior to their extinction.
The Caribbean ground sloths, the most recent survivors, lived on Cuba and Hispaniola, possibly until 1550 BCE. However, radiocarbon dating suggests an age of between 2819 and 2660 BCE for the last occurrence of Megalocnus in Cuba. They survived 5,000–6,000 years longer in the Caribbean than on the American mainland, which correlates with the later colonization of this area by humans.
Description
Ground sloths varied widely in size from under in the Caribbean ground sloths, to in the largest ground sloth genera Megatherium, Lestodon and Eremotherium. The bodies of ground sloths were generally barrel-shaped, with a broad pelvis. Like other xenarthrans, the adult teeth of ground sloths lacked enamel, with the tooth surface being composed of relatively soft orthodentine. The number of teeth in the jaws is considerably reduced in comparison to other mammals, with most ground sloths only having 5 and 4 teeth in each half of the upper and lower jaws respectively, with some ground sloths exhibiting further tooth number reduction. These teeth were rootless and were continuously growing (hypselodont), and typically have a relatively simple morphology. Some ground sloths have canine-like teeth at the front of the jaws separated from the other teeth by a gap (diastema). The skull shapes of ground sloths are highly variable. Those with narrow muzzles are likely to have had prehensile lips, while those with wider muzzles are likely to have had mobile tongues. The hands of ground sloths have ungual phalanges that indicate that they had well developed claws. In many ground sloth families (Megatheriidae, Mylodontidae, Scelidotheriidae and Nothrotheriidae), the hindfoot is inwardly rotated, meaning sole faces inwards and that the body weight was primarily borne on the fifth metatarsus and the calcaneum.
Ecology
Ground sloths are generally regarded as herbivores, with some being browsers, others grazers, and some intermediate between the two as mixed feeders (both browsing and grazing), though a number of authors have argued that some ground sloths may have been omnivores. Sloths that had longer snouts are presumed to have had greater olfactory acuity, but appear to have also had less binocular vision and poorer ability to localize sounds. A number of extinct sloth species are thought to have had hearing abilities optimized for low frequencies, perhaps related to use of infrasound for communication. Some ground sloths are suggested to have dug burrows. Their skeletal anatomy suggests that they were incapable of running, and relied on other strategies to defend against predators, though they were likely significantly more active and agile than living tree sloths. Ground sloths were likely able to adopt a bipedal stance while stationary, allowing the forelimbs to be used to grasp vegetation as well as to use their claws for defence, though whether they were capable of moving in this posture is uncertain. Some ground sloths have been suggested to be able to climb. Some authors have suggested ground sloths were largely solitary animals, like living sloths, though other authors have argued that at least some ground sloths are likely to have engaged in gregarious behaviour. Whether or not ground sloths had a slow metabolism like living xenarthrans (including living sloths) is debated.
Like living sloths, ground sloths likely only gave birth to a single offspring at a time, with likely several years between the birth of offspring. At least some ground sloths engaged in long-term parental care, with one adult (presumably female) Megalonyx found with two juveniles of different ages, with the oldest juvenile suggested to be 3–4 years old. Juvenile ground sloths may have clung to the body of their mother for some time following birth, as occurs in living tree sloths.
Evolution
The earliest unambiguous fossil evidence of ground sloths comes from the early Oligocene. Ground sloths had dispersed into the Caribbean already by 31 million years ago, as evidenced by a femur found in Puerto Rico. During the Miocene, sloths diversified, with the major families of sloths appearing during this period, with diversity waxing and waning over the course of the Miocene. Megalonychid and mylodontid sloths had migrated into North America by the Late Miocene, around 10 million years ago. At the end of the Miocene, ground sloth diversity declined, though their diversity would remain largely stable throughout the Pliocene and Pleistocene periods, up until their extinction. During the Pliocene and Pleistocene, as part of the Great American Interchange, additional lineages of sloths migrated into Central and North America. Prior to their extinction, there were over 30 living species of ground sloths across the Americas during the Late Pleistocene.
Families
Paleontologists assign more than 80 genera of ground sloths to multiple families.
Megalonychidae
The megalonychid ground sloths first appeared in the Late Eocene, about 35 million years ago, in Patagonia. Megalonychids first reached North America by island-hopping, prior to the formation of the Isthmus of Panama. Some lineages of megalonychids increased in size as time progressed. The first species of these were small and may have been partly tree-dwelling, whereas the Pliocene (about 5 to 2 million years ago) species were already approximately half the size of the huge Late Pleistocene Megalonyx jeffersonii from the last ice age. Some West Indian island species were as small as a large cat; their dwarf condition typified both tropical adaptation and their restricted island environment. This small size also enabled them a degree of arboreality.
Megalonyx, which means "giant claw", was a widespread North American genus that lived past the close of the last (Wisconsin) glaciation, when so many large mammals died out. Remains have been found as far north as Alaska and the Yukon. Ongoing excavations at Tarkio Valley in southwestern Iowa may reveal something of the familial life of Megalonyx. An adult was found in direct association with two juveniles of different ages, suggesting that adults cared for young of different generations.
The earliest known North American megalonychid, Pliometanastes protistus, lived in the southern U.S. about 9 million years ago and is believed to have been the predecessor of Megalonyx. Several species of Megalonyx have been named; in fact it has been stated that "nearly every good specimen has been described as a different species". A broader perspective on the group, accounting for age, sex, individual and geographic differences, indicates that only three species are valid (M. leptostomus, M. wheatleyi, and M. jeffersonii) in the late Pliocene and Pleistocene of North America, although work by McDonald lists five species.
Jefferson's ground sloth has a special place in modern paleontology, for Thomas Jefferson's letter on Megalonyx, read before the American Philosophical Society of Philadelphia in August 1796, marked the beginning of vertebrate paleontology in North America. When Lewis and Clark set out, Jefferson instructed Meriwether Lewis to keep an eye out for ground sloths. He was hoping they would find some living in the Western range. Megalonyx jeffersonii was appropriately named after Thomas Jefferson.
Megatheriidae
The megatheriid ground sloths are relatives of the megalonychids; these two families, along with the family Nothrotheriidae, form the infraorder Megatheria. Megatheriids appeared later in the Oligocene, some 30 million years ago, also in South America. The group includes the heavily built Megatherium (given its name 'great beast' by Georges Cuvier) and Eremotherium, which are the largest known ground sloths, thought to have had body masses of 3.5-4 tons. The skeletal structure of these ground sloths indicates that the animals were massive. Their thick bones and even thicker joints (especially those on the hind legs) gave their appendages tremendous power that, combined with their size and fearsome claws, provided a formidable defense against predators.
The earliest megatheriid in North America was Eremotherium eomigrans which arrived 2.2 million years ago, after crossing the recently formed Panamanian land bridge. With more than five tons in weight, 6 meters in length, and able to reach as high as , it was larger than an African bush elephant bull. Unlike relatives, this species retained a plesiomorphic extra claw. While other species of Eremotherium had four fingers with only two or three claws, E. eomigrans had five fingers, four of them with claws up to nearly a foot long.
Nothrotheriidae
Recently recognized, ground sloths of Nothrotheriidae are often associated with those of the Megatheriidae, and together the two form the superfamily Megatherioidea. The most prominent members of the group are the South American genus Thalassocnus, known for being aquatic, and Nothrotheriops from North America.
The last ground sloths in North America belonging to Nothrotheriops died so recently that their subfossil dung has remained undisturbed in some caves. One of the skeletons, found in a lava tube (cave) at Aden Crater, adjacent to Kilbourne Hole, New Mexico, still had skin and hair preserved, and is now at the Yale Peabody Museum. The largest samples of Nothrotheriops dung can be found in the collections of the Smithsonian Museum. Another Nothrotheriops was excavated at Shelter Cave, also in Doña Ana County, New Mexico.
Mylodontidae
The mylodontid ground sloths together with their relatives the scelidotheriids form the Mylodontoidea, the second radiation of ground sloths. The discovery of their fossils in caverns associated with human occupation led some early researchers to theorize that the early humans built corrals when they could procure a young ground sloth, to raise the animal to butchering size. However, radiocarbon dates do not support simultaneous occupation of the site by humans and sloths. Subfossil remains like coproliths, fur and skin have been discovered in some quantities. The American Museum of Natural History has exhibited a sample of Mylodon dung from Argentina with a note that reads "deposited by Theodore Roosevelt". Mylodontids are the only ground sloths confirmed to have had osteoderms embedded within their skin, though osteoderms were only present in a handful of genera and absent in many others.
The largest mylodontid is Lestodon, with an estimated mass of .
Scelidotheriidae
The ground sloth family Scelidotheriidae was demoted in 1995 to the subfamily Scelidotheriinae within Mylodontidae. Based on collagen sequence data showing that its members are more distant from other mylodontids than Choloepodidae, it was elevated back to full family status in 2019. Together with Mylodontidae, the enigmatic Pseudoprepotherium and two-toed sloths, the scelidotheriids form the superfamily Mylodontoidea. Chubutherium is an ancestral and very plesiomorphic member of this subfamily and does not belong to the main group of closely related genera, which include Scelidotherium and Catonyx.
Phylogeny
The following sloth family phylogenetic tree is based on collagen and mitochondrial DNA sequence data (see Fig. 4 of Presslee et al., 2019).
Extinction
Radiocarbon dating places the disappearance of ground sloths in what is now the United States at around 11,000 years ago. The Shasta ground sloth (Nothrotheriops shastensis) visited Rampart Cave (located on the Arizona side of the Lake Mead National Recreation Area) seasonally, leaving behind a massive stratified subfossilized dung deposit, and seemed to be flourishing from 13,000 until 11,000 BP, when the deposition suddenly stopped. Steadman et al. argue that it is no coincidence that studies have shown that ground sloths disappeared from an area a few years after the arrival of humans. Trackways preserved in New Mexico (probably dating from 10 to 15.6 thousand years ago) that appear to show a group of humans chasing or harassing three Nothrotheriops or Paramylodon ground sloths may record the scene of a hunt. The tracks are interpreted as showing seven instances of a sloth turning and rearing up on its hind legs to confront its pursuers, while the humans approach from multiple directions, possibly in an attempt to distract it.
Those who argue in favor of humans being the direct cause of the ground sloths' extinction point out that the few sloths that remain are small sloths that spend most of their time in trees, making it difficult for them to be spotted. Although these sloths were well hidden, they still would have been affected by the climate changes that others claim wiped out the ground sloths. Additionally, after the continental ground sloths disappeared, insular sloths of the Caribbean survived for approximately 6,000 years longer, which correlates with the fact that these islands were not colonized by humans until about 5500 yr BP.
It is difficult to find evidence that supports either claim on whether humans hunted the ground sloths to extinction. Removing large amounts of meat from large mammals such as the ground sloth requires no contact with the bones; tool-inflicted damage to bones is a key sign of human interaction with the animal.
Hunting of ground sloths
Kill sites
A number of kill sites are known for ground sloths in the Americas, these include Campo Laborde in the Pampas of Argentina, where an individual of Megatherium americanum was butchered at the edge of a swamp, dating to approximately 12,600 years Before Present (BP), with another potential Megatherium kill site being Arroyo Seco 2 in the same region, dating to approximately 14,782–11,142 cal yr BP. In northern Ohio, a Megalonyx jeffersoni skeleton dubbed the "Firelands Ground Sloth" has cut marks indicative of butchery, dating to 13,738 to 13,435 years BP. At the Santa Elina rockshelter in Mato Grosso Brazil, a specimen of Glossotherium is associated with hearths and stone tools, dating to 11,833–11,804 years BP. At Fell's Cave in southern Chilean Patagonia, a specimen of Mylodon with fractured and burned bones associated with human activity has been dated to approximately 12,766–12,354 years BP.
Hunting weapons
Humans are believed to have entered the New World via Beringia, a land bridge which connected Asia and North America during the last glacial maximum. Mosimann and Martin (1975) suggested the first of these nomads descended from hunting families who had acquired the skills to track down and kill large mammals. By this time, humans had developed proficient hunting weapons, including the Clovis points, which were narrow, carved stone projectiles used specifically for big game. A couple of hundred years later, the atlatl became widely used, which allowed them to throw spears with greater velocity. These inventions would have allowed hunters to put distance between them and their prey, potentially making it less dangerous to approach ground sloths.
Advantages
Certain characteristics and behavioral traits of the ground sloths made them easy targets for human hunting and provided hunter-gatherers with strong incentives to hunt these large mammals.
Ground sloths often fed in open fields. Recent studies have attempted to discover the diet of ground sloths through fossils of their dung. Analysis of these coproliths have found that ground sloths often ate the foliage of trees, hard grasses, shrubs, and yucca; these plants were located in areas that would have exposed them, making them susceptible to human predation. Ground sloths were not only easy to spot, but had never interacted with humans before, so would not have known how to react to them. Additionally, these large mammals waddled on their hind legs and front knuckles, keeping their claws turned in. Their movement and massive build (some weighed up to ) imply they were relatively slow mammals.
These reasonable after-the-fact inferences from the evidence might explain why ground sloths would have been easy prey for hunters, but are not certain.
Difficulties
While ground sloths would have been relatively easy to spot and approach, big game hunters' weapons would have been useless from farther than away. It would have been difficult to take down a ground sloth with a spear-thrower and would have required extensive knowledge of the species. Additionally, the ground sloths' already thick hide was fortified by osteoderms, making it difficult to penetrate.
Since ground sloths thrived in an environment filled with large predators, they evidently would have been able to also defend themselves against human predation, so there is no reason to expect that they would have been "easy pickings". When feeding, they had enough strength to use their long, sharp claws to tear apart tree branches; presumably their strength and formidable claws would be dangerous for hunters that attempted to attack them at close quarters. But fossilized evidence of humans hunting on ground sloth in White Sands National Park suggests that the slow-moving giant sloths were likely easy prey for early humans possibly hurling spears.
| Biology and health sciences | Xenarthra | Animals |
276589 | https://en.wikipedia.org/wiki/Rusty-spotted%20cat | Rusty-spotted cat | The rusty-spotted cat (Prionailurus rubiginosus) is one of the cat family's smallest members, of which historical records are known only from India and Sri Lanka. In 2012, it was also recorded in the western Terai of Nepal. Since 2016, the global wild population is listed as Near Threatened on the IUCN Red List as it is fragmented and affected by loss and destruction of its prime habitat, deciduous forests.
Taxonomy
Felis rubiginosa was the scientific name used by Isidore Geoffroy Saint-Hilaire in 1831 for a rusty-spotted cat specimen from Pondicherry, India. Prionailurus was proposed by Nikolai Severtzov in 1858 as a generic name. Prionailurus rubiginosus phillipsi was proposed by Reginald Innes Pocock in 1939 who described a specimen from Central Province, Sri Lanka and subordinated both to the genus Prionailurus.
Phylogeny
Phylogenetic analysis of the nuclear DNA in tissue samples from all Felidae species revealed that the evolutionary radiation of the Felidae began in Asia in the Miocene around . Analysis of mitochondrial DNA of all Felidae species indicates a radiation at around .
The Prionailurus species are estimated to have had a common ancestor between , and . The rusty-spotted cat possibly genetically diverged from this ancestor between . Both models agree in the rusty-spotted cat having been the first cat of this lineage that diverged, followed by the flat-headed cat (P. planiceps) and the fishing cat (P. viverrinus). The following cladogram shows the phylogenetic relationships of the rusty-spotted cat as derived through analysis of nuclear DNA:
Characteristics
The rusty-spotted cat has a short reddish grey fur over most of the body with rusty spots on the back and flanks. Four blackish lines run over the eyes, and two of them extend over the neck. Six dark streaks are on each side of the head, extending over the cheeks and forehead. Its chin, throat, inner side of the limbs and belly are whitish with tiny brownish spots. It has a rusty band on the chest. Its paws and tail are uniform reddish grey.
It is the smallest wild cat in Asia and rivals the black-footed cat as the world's smallest wild cat. It is in length, with a tail, and weighs only . The bushy tail is about half the length of the body.
Distribution and habitat
The distribution of the rusty-spotted cat is relatively restricted. It occurs mainly in moist and dry deciduous forests as well as scrub and grassland, but is likely absent in evergreen forest. It prefers dense vegetation and rocky areas.
In India, it was long thought to be confined to the south, but records have established that it occurs over much of the country. It was observed in eastern Gujarat's Gir National Park, in Maharashtra's Tadoba-Andhari Tiger Reserve and along India's Eastern Ghats. Camera trapping revealed its presence in Pilibhit Tiger Reserve in the Indian Terai and in Nagzira Wildlife Sanctuary in Maharashtra.
In western Maharashtra, the rusty-spotted cat is breeding in a human dominated agricultural landscape, where rodent densities are high. In December 2014 and in April 2015, it was photographed by camera traps in Kalesar National Park, Haryana. It was also recorded by camera traps in Mirzapur Forest Division of Uttar Pradesh in 2018.
In March 2012, a rusty-spotted cat was photographed in Bardia National Park for the first time, and in March 2016 also in Shuklaphanta National Park, both in Nepal.
In Sri Lanka, there are a few records in montane and lowland rainforest. There are two distinct populations, one in the dry zone and the other in the wet zone.
In 2016, it was recorded for the first time in Horton Plains National Park at elevations of .
Ecology and behaviour
Very little is known about the ecology and behaviour of the rusty-spotted cat in the wild. Captive ones are mostly nocturnal but also briefly active during the day. Most wild ones were also recorded after dark. At Horton Plain National Park in Sri Lanka, they were mostly recorded between sunset and sunrise, with limited daytime activity. Several individuals were observed hiding in trees and in caves.
It feeds mainly on rodents and birds, but also preys on lizards, frogs, and insects. It hunts primarily on the ground, making rapid, darting movements to catch its prey. It apparently ventures into trees to escape larger predators. Captive females and males both scent-mark their home range by spraying urine.
Reproduction
The female's oestrus lasts five days, and mating is unusually brief. Since the female is likely to be vulnerable during this period, its brevity may be an adaptation to help it avoid larger predators. She prepares a den in a secluded location, and after a gestation of 65–70 days gives birth to one or two kittens. At birth, the kittens weigh just , and are marked with rows of black spots. They reach sexual maturity at around 68 weeks, by which time they have developed the distinctive adult coat pattern of rusty blotches. Rusty-spotted cats have lived for twelve years in captivity, but their lifespan in the wild is unknown.
Threats
Habitat loss and the spread of cultivation are serious problems for wildlife in both India and Sri Lanka. Although there are several records of the rusty-spotted cat in cultivated and settled areas, it is not known to what degree these populations are able to persist in such areas. There have been occasional reports of rusty-spotted cat skins in trade. In some areas, it is hunted for food or as livestock pest.
Conservation
The Indian population is listed on CITES Appendix I. The Sri Lankan population is included on CITES Appendix II. The species is fully protected over most of its range, with hunting and trade banned in India and Sri Lanka.
As of 2010, the captive population of P. r. phillipsi comprised 56 individuals in eight institutions, of which 11 individuals were kept in the Colombo Zoo in Sri Lanka and 45 individuals in seven European zoos.
| Biology and health sciences | Felines | Animals |
276610 | https://en.wikipedia.org/wiki/Fishing%20cat | Fishing cat | The fishing cat (Prionailurus viverrinus) is a medium-sized wild cat of South and Southeast Asia. It has a deep yellowish-grey fur with black lines and spots. Adults have a head-to-body length of , with a long tail. Males are larger than females, weighing , while females average . It lives mostly in the vicinity of wetlands, along rivers, streams, oxbow lakes, in swamps and mangroves where it preys mostly on fish. Other prey items include birds, insects, small rodents, molluscs, reptiles including snakes, amphibians and carrion of cattle. The fishing cat is thought to be primarily nocturnal. It is a good swimmer and can swim long distances, even underwater.
The fishing cat has been listed as a vulnerable species on the IUCN Red List since 2016, as the global population is thought to have declined by about 30% in the past three fishing cat generations during the period 2010–2015. The destruction of wetlands and killing by local people are the major threats throughout its range. It is the state animal of West Bengal.
Taxonomy
Felis viverrinus was proposed by Edward Turner Bennett in 1833, who described a cat skin sent from India by Josiah Marshall Heath. The genus name Prionailurus was proposed by Nikolai Severtzov in 1858 for spotted wild cats native to Asia. A subspecies Felis viverrinus rhizophoreus was proposed in 1936 by Henri Jacob Victor Sody, who described a specimen from the north coast of West Java that had a slightly shorter skull than fishing cat specimens from Thailand. There is evidence that the nominate taxon and the Javan fishing cat are distinguishable by skull morphometrics.
Phylogeny
Phylogenetic analysis of the nuclear DNA in tissue samples from all Felidae species revealed that the evolutionary radiation of the Felidae began in Asia in the Miocene around . Analysis of mitochondrial DNA of all Felidae species indicates a radiation at around .
The Prionailurus species are estimated to have had a common ancestor between , and .
Both models agree in the rusty-spotted cat (P. rubiginosus) having been the first cat of the Prionailurus lineage that genetically diverged, followed by the flat-headed cat (P. planiceps) and then the fishing cat. It is estimated to have diverged together with the leopard cat (P. bengalensis) between and .
The following cladogram shows the phylogenetic relationships of the fishing cat as derived through analysis of nuclear DNA:
Characteristics
The fishing cat has a deep yellowish-grey fur with black lines and spots. Two stripes are on the cheeks, and two are above the eyes running to the neck with broken lines on the forehead. It has two rows of spots around the throat. The spots on the shoulder are longitudinal, and those on the sides, limbs and tail are roundish. The background colour of its fur varies between individuals from yellowish tawny to ashy grey, and the size of the stripes ranges from narrow to broad. The fur on the belly is lighter than on the back and sides. The short and rounded ears are set low on the head, and the back of the ears bear a white spot. The tail is short, less than half the length of head and body, and with a few black rings at the end. The short dense layered fur is thought to be an aquatic adaptation providing a water barrier and thermal insulation, while another layer of protruding long guard hairs provides its pattern and glossy sheen.
The fishing cat is the largest cat of the Prionailurus. It is stocky and muscular with a head-to-body length ranging from , a tail length of and medium to short legs. Females weigh and males , showing pronounced sexual dimorphism. Its skull is elongated, with a basal length of and a post-orbital width of . Its tongue is about long and has large, cylindrical papillae near the front.
Fishing cat paws are partly webbed, and the claws are incompletely sheathed, only becoming partially covered when fully retracted. It is about twice the size of a domestic cat.
Distribution and habitat
The fishing cat is broadly but discontinuously distributed in South and Southeast Asia. It predominantly inhabits densely vegetated wetlands around slow-moving bodies of water like swamps and marshes. These include low-salinity bodies such as oxbow lakes, and high-salinity ones such as tidal creeks and mangrove forests. Along these bodies of water, it conceals itself within the thick cover of forests, scrublands, reed beds, and grasslands. Most records are from lowland areas.
In Pakistan's Sindh Province, the fishing cat was recorded in the Chotiari Dam area in 2012. In the Nepal Terai, it has been recorded in Shuklaphanta, Bardia, Chitwan and Parsa National Parks and in Koshi Tappu Wildlife Reserve. In India, its presence has been documented in Ranthambore National Park, in Pilibhit, Dudhwa and Valmiki Tiger Reserves, in Sur Sarovar Bird Sanctuary, outside protected areas in West Bengal, in Lothian Island Wildlife Sanctuary in the Sundarbans, in Odisha's Bhitarkanika Wildlife Sanctuary and coastal districts outside protected areas, in Andhra Pradesh's Coringa Wildlife Sanctuary, Krishna Wildlife Sanctuary and adjoining reserve forests. Reports in Bangladeshi newspapers indicate that fishing cats live in most divisions of Bangladesh. In Sri Lanka, it has been recorded in tea estates, Maduru Oya National Park and multiple localities in coastal to hilly regions.
In Myanmar, it was recorded in the Ayeyarwady Delta in 2016 and 2018. In Thailand, its presence has been documented in Khao Sam Roi Yot National Park, Thale Noi Non-Hunting Area and in Kaeng Krachan National Park. Between 2007 and 2016, it was also recorded near wetlands outside protected areas in Phitsanulok Province, Bang Khun Thian District, Samut Sakhon, Phetchaburi and Songkhla Provinces, and near a mangrove site in Pattani. In Cambodia, a single fishing cat was photographed by a camera trap in Kulen Promtep Wildlife Sanctuary in March 2003. In 2015, it was also recorded in Peam Krasop Wildlife Sanctuary and Ream National Park. The island of Java constitutes the southern limit of the fishing cat's range, but by the 1990s fishing cats were thought to be scarce and restricted to tidal forests with sandy or muddy shores, older mangrove stands, and abandoned mangrove plantation areas with fishponds. There are no confirmed records in Peninsular Malaysia, Vietnam and Laos.
Behaviour and ecology
The fishing cat lives among dense vegetation near water and is thought to be primarily nocturnal. It is known to be a proficient long-distance and underwater swimmer. Adult males and females without dependent young are solitary. Females have been reported to range over areas of , while males range over . It has been observed resting in thick grassy habitats, often near a water body but sometimes far away from them. Adults have been recorded to vocalise "chuckling" sounds.
It marks its home range using cheek-rubbing, head rubbing, chin rubbing, neck rubbing and urine-spraying to leave scent marks; it also sharpens its claws and displays flehmen. The pungent smell of fishing cat urine markings is due to 3-Mercapto-3-methylbutan-1-ol, a breakdown product of felinine.
Fishing cat feces collected in India's Keoladeo National Park revealed that fish comprises about three-quarter of its diet, with the remainder consisting of birds, small rodents and insects; molluscs, reptiles, amphibians, carrion of cattle and grass supplement its diet. Its main prey in the Godavari River delta is fish, including flathead grey mullet, green chromide and Mozambique tilapia, which comprised three fifths of its diet, whereas rodents and crabs made up the remainder of the diet. The diet make-up remained relatively constant throughout the year. Fishing cats have been observed while hunting along the edges of watercourses, grabbing prey from the water, and sometimes diving into the water to catch prey further from the banks. It prefers hunting in shallow water and spends about half the time lying in wait for prey to approach.
Reproduction and development
Wild fishing cats most likely mate during January and February; most kittens in the wild were observed in March and April. However, fishing cats may mate as late as June. In captivity, the gestation period lasts 63–70 days; females give birth to an average of two to three kittens; the litter size can be as small as one to as large as four. Kittens weigh around at birth and are able to actively move around by the age of one month. They begin to play in water and to take solid food when about two months old, but are not fully weaned until six months old. They reach full adult size when about eight and a half months old, acquire their adult canine teeth by 11 months and are sexually mature when approximately 15 months old. They live up to 10 years in captivity. The generation length of the fishing cat is five years.
Health
Fishing cats are susceptible to carnivore protoparvovirus, a disease known to kill them. This disease significantly damages the kidney, spleen and gastrointestinal tract of the body. The fishing cat is also vulnerable to diseases and medical conditions such as feline hemoplasmas, transitional cell carcinoma and canine distemper virus. One captive individual was even recorded with chlamydiota. Additionally, in a 2012 case study, Toxocara cati was reported as the cause of death of young captive fishing cats.
Threats
Since 2016, the fishing cat is listed as a vulnerable species on the IUCN Red List, as the global population is thought to have declined by around 30% in the years 2010–2015; the destruction of wetlands and killings by local people are major threats to the fishing cat.
The destruction of wetlands includes increased pollution or conversion for agricultural use and human settlements. The conversion of mangrove forests to commercial aquaculture ponds is a major threat in Andhra Pradesh, and in some places fishing cats are killed. Over-exploitation of local fish stocks and retaliatory killing by local people are also significant threats. The fishing cat's habitat in India is predominantly marshlands, which are subject to agricultural usage under the country's laws, resulting in human–wildlife conflict. Coastal wetlands are an important habitat for the fishing cat in Thailand and Cambodia, however, estimates predict that only 6% of wetlands remain undisturbed.
In West Bengal's Howrah district, 27 dead fishing cats were recorded between April 2010 and May 2011, and in Sagar Island, the fishing cat was possibly extirpated by local people for reasons unknown. Deaths are more frequent in the dry season, when people use wetlands more frequently. At least 30 fishing cats were killed by local people in Bangladesh between January 2010 and March 2013. They are often killed because they are mistaken as tiger cubs. In Thailand, 31% of radio-collared fishing cats were killed by local people between 2012 and 2015. Fish farmers in Koshi Tappu Wildlife Reserve have killed fishing cats in retaliation for perceived loss of fish. Roadkills are a major mortality factor in Odisha. The fishing cat is possibly extinct in coastal Kerala, but it is doubtful whether it ever occurred there at all.
Conservation
The fishing cat is included on CITES Appendix II and protected by national legislation over most of its range. Hunting is prohibited in Bangladesh, China, India, Indonesia, Myanmar, Nepal, Pakistan, Sri Lanka, and Thailand. Hunting regulations apply in Laos. In Bhutan, Malaysia, and Vietnam, it is not protected outside protected areas, and no information is known about its legal protections in Cambodia. It is the state animal of the Indian state of West Bengal. Its survival depends on protection of wetlands, prevention of indiscriminate trapping, snaring and poisoning.
In areas where habitat degradation is a major concern, such as coastal Andhra Pradesh, NGOs are working to slow habitat conversion in collaboration with local villagers. Part of this work involves creating alternative livelihood programs that allow villagers to earn money without damaging natural habitats. A Fishing Cat Conservation Alliance provides an umbrella for the cooperation of national fishing cat conservation groups, which began with the establishment of India's in 2010.
Fishing cat captive breeding programmes have been established by the European Association of Zoos and Aquaria and the American Association of Zoos and Aquariums. All the fishing cats kept in zoos around the world are listed in the International Studbook of the World Association of Zoos and Aquariums. Zoos in Thailand house around 30 individuals; birth rates are not particularly high. They have been placed in captivity as an "insurance population" due to their vulnerable status in the wild.
| Biology and health sciences | Felines | Animals |
276745 | https://en.wikipedia.org/wiki/Homo%20ergaster | Homo ergaster | Homo ergaster is an extinct species or subspecies of archaic humans who lived in Africa in the Early Pleistocene. Whether H. ergaster constitutes a species of its own or should be subsumed into H. erectus is an ongoing and unresolved dispute within palaeoanthropology. Proponents of synonymisation typically designate H. ergaster as "African Homo erectus" or "Homo erectus ergaster". The name Homo ergaster roughly translates to "working man", a reference to the more advanced tools used by the species in comparison to those of their ancestors. The fossil range of H. ergaster mainly covers the period of 1.7 to 1.4 million years ago, though a broader time range is possible. Though fossils are known from across East and Southern Africa, most H. ergaster fossils have been found along the shores of Lake Turkana in Kenya. There are later African fossils, some younger than 1 million years ago, that indicate long-term anatomical continuity, though it is unclear if they can be formally regarded as H. ergaster specimens. As a chronospecies, H. ergaster may have persisted to as late as 600,000 years ago, when new lineages of Homo arose in Africa.
Those who believe H. ergaster should be subsumed into H. erectus consider there to be too little difference between the two to separate them into distinct species. Proponents of keeping the two species as distinct cite morphological differences between the African fossils and H. erectus fossils from Asia, as well as early Homo evolution being more complex than what is implied by subsuming species such as H. ergaster into H. erectus. Additionally, morphological differences between the specimens commonly seen as constituting H. ergaster might suggest that H. ergaster itself does not represent a cohesive species. Regardless of their most correct classification, H. ergaster exhibit primitive versions of traits later expressed in H. erectus and are thus likely the direct ancestors of later H. erectus populations in Asia. Additionally, H. ergaster is likely ancestral to later hominins in Europe and Africa, such as modern humans and Neanderthals.
Several features distinguish H. ergaster from australopithecines as well as earlier and more basal species of Homo, such as H. habilis. Among these features are their larger body mass, relatively long legs, obligate bipedalism, relatively small jaws and teeth (indicating a major change in diet) as well as body proportions and inferred lifestyles more similar to modern humans than to earlier and contemporary hominins. With these features in mind, some researchers view H. ergaster as being the earliest true representative of the genus Homo.
H. ergaster lived on the savannah in Africa, a unique environment with challenges that would have resulted in the need for many new and distinct behaviours. Earlier Homo probably used counter-attack tactics, like modern primates, to keep predators away. By the time of H. ergaster, this behaviour had probably resulted in the development of true hunter-gatherer behaviour, a first among primates. H. ergaster was an apex predator. Further behaviours that might first have arisen in H. ergaster include male-female divisions of foraging and true monogamous pair bonds. H. ergaster also marks the appearance of more advanced tools of the Acheulean industry, including the earliest known hand axes. Though undisputed evidence is missing, H. ergaster might also have been the earliest hominin to master control of fire.
Taxonomy
Research history
The systematics and taxonomy of Homo in the Early to Middle Pleistocene is one of the most disputed areas of palaeoanthropology. In early palaeoanthropology and well into the twentieth century, it was generally assumed that Homo sapiens was the end result of gradual modifications within a single lineage of hominin evolution. As the perceived transitional form between early hominins and modern humans, H. erectus, originally assigned to contain archaic human fossils in Asia, came to encompass a wide range of fossils covering a large span of time (almost the entire temporal range of Homo). Since the late twentieth century, the diversity within H. erectus has led some to question what exactly defines the species and what it should encompass. Some researchers, such as palaeoanthropologist Ian Tattersall in 2013, have questioned H. erectus since it contains an "unwieldly" number of fossils with "substantially differing morphologies".
In the 1970s, palaeoanthropologists Richard Leakey and Alan Walker described a series of hominin fossils from Kenyan fossil localities on the eastern shore of Lake Turkana. The most notable finds were two partial skulls; KNM ER 3733 and KNM ER 3883, found at Koobi Fora. Leakey and Walker assigned these skulls to H. erectus, noting that their brain volumes (848 and 803 cc respectively) compared well to the far younger type specimen of H. erectus (950 cc). Another significant fossil was a fossil mandible recovered at Ileret and described by Leakey with the designation KNM ER 992 in 1972 as "Homo of indeterminate species".
In 1975, palaeoanthropologists Colin Groves and Vratislav Mazák designated KNM ER 992 as the holotype specimen of a distinct species, which they dubbed Homo ergaster. The name (ergaster being derived from the Ancient Greek ἐργαστήρ, ergastḗr, 'workman') roughly translates to "working man" or "workman". Groves and Mazák also included many of the Koobi Fora fossils, such as KNM ER 803 (a partial skeleton and some isolated teeth) in their designation of the species, but did not provide any comparison with the Asian fossil record of H. erectus in their diagnosis, inadvertently causing some of the later taxonomic confusion in regards to the species.
A nearly complete fossil, interpreted as a young male (though the sex is actually undetermined), was discovered at the western shore of Lake Turkana in 1984 by Kenyan archaeologist Kamoya Kimeu. The fossils were described by Leakey and Walker, alongside paleoanthropologists Frank Brown and John Harris, in 1985 as KNM-WT 15000 (nicknamed "Turkana Boy"). They interpreted the fossil, consisting of a nearly complete skeleton, as representing H. erectus. Turkana Boy was the first discovered comprehensively preserved specimen of H. ergaster/erectus found and constitutes an important fossil in establishing the differences and similarities between early Homo and modern humans. Turkana Boy was placed in H. ergaster by paleoanthropologist Bernard Wood in 1992, and is today, alongside other fossils in Africa previously designated as H. erectus, commonly seen as a representative of H. ergaster by those who support H. ergaster as a distinct species.
Classification
H. ergaster is easily distinguished from earlier and more basal species of Homo, notably H. habilis and H. rudolfensis, by a number of features that align them, and their inferred lifestyle, more closely to modern humans than to earlier and contemporary hominins. As compared to their relatives, H. ergaster had body proportions more similar to later members of the genus Homo, notably relatively long legs which would have made them obligately bipedal. The teeth and jaws of H. ergaster are also relatively smaller than those of H. habilis and H. rudolfensis, indicating a major change in diet. In 1999, palaeoanthropologists Bernard Wood and Mark Collard argued that the conventional criteria for assigning species to the genus Homo were flawed and that early and basal species, such as H. habilis and H. rudolfensis, might appropriately be reclassified as ancestral australopithecines. In their view, the true earliest representative of Homo was H. ergaster.
Since its description as a separate species in 1975, the classification of the fossils referred to H. ergaster has been in dispute. H. ergaster was immediately dismissed by Leakey and Walker and many influential researchers, such as palaeoanthropologist G. Philip Rightmire, who wrote an extensive treatise on H. erectus in 1990, continued to prefer a more inclusive and comprehensive H. erectus. Overall, there is no doubt that the group of fossils composing H. erectus and H. ergaster represent the fossils of a more or less cohesive subset of closely related archaic humans. The question is instead whether these fossils represent a radiation of different species or the radiation of a single, highly variable and diverse, species over the course of almost two million years. This long-running debate remains unresolved, with researchers typically using the terms H. erectus s.s. (sensu stricto) to refer to H. erectus fossils in Asia and the term H. erectus s.l. (sensu lato) to refer to fossils of other species that may or may not be included in H. erectus, such as H. ergaster, H. antecessor and H. heidelbergensis.
For obvious reasons, H. ergaster shares many features with H. erectus, such as large forward-projecting jaws, large brow ridges and a receding forehead. Many of the features of H. ergaster are clearly more primitive versions of features later expressed in H. erectus, which somewhat obscures the differences between the two. There are subtle, potentially significant, differences between the East African and East Asian fossils. Among these are the somewhat higher-domed and thinner-walled skulls of H. ergaster, and the even more massive brow ridges and faces of Asian H. erectus.
The question is made more difficult since it regards how much intraspecific variation can be exhibited in a single species before it needs to be split into more, a question that in and of itself does not have a clear-cut answer. A 2008 analysis by anthropologist Karen L. Baab, examining fossils of various H. erectus subspecies, and including fossils attributed to H. ergaster, found that the intraspecific variation within H. erectus was greater than expected for a single species when compared to modern humans and chimpanzees, but fell well within the variation expected for a species when compared to gorillas, and even well within the range expected for a single subspecies when compared to orangutans (though this is partly due to the great sexual dimorphism exhibited in gorillas and orangutans). Baab concluded that H. erectus s.l. was either a single but variable species, several subspecies divided by time and geography or several geographically dispersed but closely related species. In 2015, paleoanthropologists David Strait, Frederick Grine and John Fleagle listed H. ergaster as one of the seven "widely recognized" species of Homo, alongside H. habilis, H. rudolfensis, H. erectus, H. heidelbergensis, H. neanderthalensis and H. sapiens, noting that other species, such as H. floresiensis and H. antecessor, were less widely recognised or more poorly known.
Variation in the fossil material
Comparing various African fossils attributed to H. erectus or H. ergaster to Asian fossils, notably the type specimen of H. erectus, in 2013, Ian Tattersall concluded that referring to the African material as H. ergaster rather than "African H. erectus" was a "considerable improvement" as there were many autapomorphies distinguishing the material of the two continents from one another. Tattersall believes it to be appropriate to use the designation H. erectus only for eastern Asian fossils, disregarding its previous use as the name for an adaptive grade of human fossils from throughout Africa and Eurasia. Though Tattersall concluded that the H. ergaster material represents the fossils of a single clade of Homo, he also found there to be considerable diversity within this clade; the KNM ER 992 mandible accorded well with other fossil mandibles from the region, such as OH 22 from Olduvai and KNM ER 3724 from Koobi Fora, but did not necessarily match with cranial material, such as KNM ER 3733 and KNM ER 3883 (since neither preserves the jaw), nor with the mandible preserved in Turkana Boy, which has markedly different dentition.
The most "iconic" fossil of H. ergaster is the KNM ER 3733 skull, which is sharply distinguished from Asian H. erectus by a number of characteristics, including that the brow ridges project forward as well as upward and arc separately over each orbit and the braincase being quite tall compared to its width, with its side walls curving. KNM ER 3733 can be distinguished from KNM ER 3883 by a number of features as well, notably in that the margins of KNM ER 3883's brow ridges are very thickened and protrude outwards but slightly downwards rather than upwards. Both skulls can be distinguished from the skull of Turkana Boy, which possesses only slightly substantial thickenings of the superior orbital margins, lacking the more vertical thickening of KNM ER 3883 and the aggressive protrusion of KNM ER 3733. In addition to this, the facial structure of Turkana Boy is narrower and longer than that of the other skulls, with a higher nasal aperture and likely a flatter profile of the upper face. It is possible that these differences can be accounted for through Turkana Boy being a subadult, 7 to 12 years old. Furthermore, KNM ER 3733 is presumed to have been the skull of a female (whereas Turkana Boy is traditionally interpreted as male), which means that sexual dimorphism may account for some of the differences.
The differences between Turkana Boy's skull and KNM ER 3733 and KNM ER 3883, as well as the differences in dentition between Turkana Boy and KNM ER 992 have been interpreted by some, such as paleoanthropologist Jeffrey H. Schwartz, as suggesting that Turkana Boy and the rest of the H. ergaster material does not represent the same taxon. Schwartz also noted none of the fossils seemed to represent H. erectus either, which he believed was in need of significant revision. In 2000, French palaeoanthropologist Valéry Zeitoun suggested that KNM ER 3733 and KNM ER 3883 should be referred to two separate species, which she dubbed H. kenyaensis (type specimen KNM ER 3733) and H. okotensis (type specimen KNM ER 3883), but these designations have found little acceptance.
Evolutionary history
Evolution and temporal range
Although frequently assumed to have originated in East Africa, the origins of H. ergaster are obscured by the fact that the species marks a radical departure from earlier species of Homo and Australopithecus in its long limbs, height and modern body proportions. Though a large number of Pleistocene tools have been found in East Africa, it can not be fully ascertained that H. ergaster originated there without further fossil discoveries. It is assumed that H. ergaster evolved from earlier species of Homo, probably H. habilis. Though populations of H. ergaster outside of Africa have been inferred based on the geographical distribution of their descendants and tools matching those in East Africa, fossils of the species are mainly from East Africa in the time range of 1.8 to 1.7 million years ago. Most fossils have been recovered from around the shores of Lake Turkana in Kenya.
The oldest known specimen of H. erectus s.l. in Africa (i.e. H. ergaster) is DNH 134, a skull recovered in the Drimolen Palaeocave System in South Africa, dated to 2.04 to 1.95 million years ago. The skull is also the oldest known H. erectus s.l. specimen overall, showing clear similarities to KNM ER 3733, and demonstrates that early H. ergaster coexisted with other hominins such as Paranthropus robustus and Australopithecus sediba.
There are also younger specimens of H. ergaster; notably, Turkana Boy is dated to about 1.56 million years ago. A handful of even younger African skulls make the case for long-term anatomical continuity, though it is unclear if they can appropriately be formally regarded as H. ergaster specimens; the "Olduvai Hominid 9" skull from Olduvai Gorge is dated to about 1.2 to 1.1 million years ago and there are also skulls from Buia (near the coast of Eritrea, dated to ~1 million years old), the Bouri Formation in Ethiopia (dated to between 1 million and 780,000 years old) and a fragmentary skull from Olorgesailie in Kenya (dated to between 970,000 and 900,000 years ago). The Olduvai skull is similar to Asian H. erectus in its massive brow ridge, but the others only show minor differences to earlier H. ergaster skulls.
The H. erectus in Asia, as well as later hominins in Europe (i. e. H. heidelbergensis and H. neanderthalensis) and Africa (H. sapiens) are all probably lineages descended from H. ergaster. Because H. ergaster is thought to have been ancestral to these later Homo, it might have persisted in Africa until around 600,000 years ago, when brain size increased rapidly and H. heidelbergensis emerged.
Expansion out of Africa
Traditionally, H. erectus was seen as the hominin that first left Africa to colonise Europe and Asia. If H. ergaster is distinct from H. erectus, this role would apply to H. ergaster instead. Very little concrete information is known on when and which Homo first appeared in Europe and Asia, since Early Pleistocene fossil hominins are scarce on both continents, and that it would have been H. ergaster (or "early H. erectus") that expanded, as well as the particular manner in which they did, remains conjecture. The presence of H. erectus fossils in East Asia means that a human species, most likely H. ergaster, had left Africa before 1 million years ago, the assumption historically having been that they first migrated out of Africa around 1.9 to 1.7 million years ago. Discoveries in Georgia and China push the latest possible date further back, before 2 million years ago, also casting doubt on the idea that H. ergaster was the first hominin to leave Africa.
The main reason for leaving Africa is likely to have been an increasing population periodically outgrowing their resource base, with splintering groups moving to establishing themselves in neighboring, empty territories over time. The physiology and improved technology of H. ergaster might have allowed them to travel to and colonise territories that no one had ever occupied before. It is unclear if H. ergaster was truly uniquely capable of expanding outside Africa; australopithecines had likely colonised savannah grasslands throughout Africa by 3 million years ago and there are no clear reasons as to why they would not have been able to expand into the grasslands of Asia before H. ergaster.
The general assumption is that hominins migrated out of the continent either across the southern end of the Red Sea or along the Nile Valley, but there are no fossil hominins known from either region in the Early Pleistocene. The earliest Homo fossils outside Africa are the Dmanisi skulls from Georgia (dated to 1.77–1.85 million years old, representing either early H. ergaster or a new taxon, H. georgicus), three incisors from Ubeidiya in Israel (about 1.4 to 1 million years old) and the fossils of Java Man (H. erectus erectus, more than five thousand miles away). The dating of key Asian H. erectus specimens (including Java Man) is not entirely certain, but they are all likely to be 1.5 million years old or younger. Ubeidiya is also the oldest firmly confirmed site of Acheulean tools (one of the tool industries associated with H. ergaster) outside Africa, the tools recovered there closely resembling older tools discovered in East Africa.
The earliest fossil evidence of Homo in Asia are the aforementioned Dmanisi skulls, which share many traits with H. ergaster in Africa, suggesting that H. ergaster might have expanded out of Africa as early as 1.7–1.9 million years ago. In addition to H. ergaster-like traits, the Dmanisi skulls possess a wide assortment of other traits, some of which are similar to traits in earlier hominins such as H. habilis, and the site notably lacks preserved hand axes (otherwise characteristic of H. ergaster), which means that hominins might have spread out of Africa even earlier than H. ergaster. The skull D2700 (Dmanisi skull 3) in particular resembles H. habilis in the small volume of its braincase (600 cc), the form of the middle and upper face and the lack of an external nose. The mixture of skulls at Dmanisi suggests that the definition of H. ergaster (or H. erectus) might most appropriately be expanded to contain fossils that would otherwise be assigned to H. habilis or that two separate species of archaic humans left Africa early on. In addition to the Dmanisi fossils, stone tools manufactured by hominins have been discovered on the Loess Plateau in China and dated to 2.12 million years old, meaning that hominins must have left Africa before that time.
An alternative hypothesis historically has been that Homo evolved in Asia from earlier ancestors that had migrated there from Africa, and then expanded back into Europe, where it gave rise to H. sapiens. This view was notably held by Eugène Dubois, who first described H. erectus fossils in the 19th century and considered the fossils of Java Man, at the time undeniably the earliest known hominin fossils, as proof of the hypothesis. Though the discovery of australopithecines and earlier Homo in Africa meant that Homo itself did not originate in Asia, the idea that H. erectus (or H. ergaster) in particular did, and then expanded back into Africa, has occasionally resurfaced. Various fossil discoveries have been used to support it through the years, including a massive set of jaws from Indonesia which were perceived to be similar to those of australopithecines and dubbed Meganthropus (now believed to be an unrelated hominid ape). The discovery of H. floresiensis in 2003, which preserved primitive foot and wrist anatomy reminiscent of that of H. habilis and Australopithecus again led to suggestions of pre-erectus hominins in Asia, though there are no known comparable foot or wrist bones from H. erectus which makes comparisons impossible. The idea that H. ergaster/H. erectus first evolved in Asia before expanding back into Africa was substantially weakened by the dating of the DNH 134 skull as approximately 2 million years old, predating all other known H. ergaster/H. erectus fossils.
Anatomy
Build and appearance
The only well-preserved post-cranial remains of H. ergaster come from the Turkana Boy fossil. Unlike the australopithecines, Turkana Boy's arms were not longer relative to their legs than the arms of living people and the cone-shaped torso of their ancestors had evolved into a more barrel-shaped chest over narrow hips, another similarity to modern humans. The tibia (shin bone) of Turkana Boy is relatively longer than the same bone in modern humans, potentially meaning that there was more bend in the knee when walking. The slim and long build of Turkana Boy may be explained by H. ergaster living in hot and arid, seasonal environments. Through thinning of the body, body volume decreases faster than skin area and greater skin area means more effective heat dissipation.
H. ergaster individuals were significantly taller than their ancestors. Whereas Lucy, a famous Australopithecus fossil, would only have been about tall at her death, Turkana Boy was about tall and would probably have reached or more if he had survived to adulthood. Adult H. ergaster are believed to have ranged in size from about tall.
Because of being adapted to a hot and arid climate, H. ergaster might also have been the earliest human species to have nearly hairless and naked skin. If instead H. ergaster had an ape-like covering of body hair, sweating (the primary means through which modern humans prevent their brains and bodies from overheating) would not have been as efficient. Though sweating is the generally accepted explanation for hairlessness, other proposed explanations include a reduction of parasite load and sexual selection. It is doubtful if australopithecines and earlier Homo were sufficiently mobile to make hair loss an advantageous trait, whereas H. ergaster was clearly adapted for long-distance travel and noted for inhabiting lower altitudes (and open, hot savannah environments) than their ancestors. Australopithecines typically inhabited colder and higher altitudes 1,000–1,600 m (3,300–5,200 ft), where nighttime temperatures would have gotten significantly colder and insulating body hair may have been required.
Alternatively and despite this, the loss of body hair could have occurred significantly earlier than H. ergaster. Though skin impressions are unknown in any extinct hominin, it is possible that human ancestors were already losing their body hair around 3 million years ago. Human ancestors acquired pubic lice from gorillas about 3 million years ago, and speciation of human from gorilla pubic lice was potentially only possible because human ancestors had lost most of their body hair by this early date. It is also possible that the loss of body hair occurred at a significantly later date. Genetic analysis suggests that high activity in the melanocortin 1 receptor, which produces dark skin, dates back to about 1.2 million years ago. This could indicate the evolution of hairlessness around this time, as a lack of body hair would have left the skin exposed to harmful UV radiation.
Skull and face
Differences to modern humans would have been readily apparent in the face and skull of H. ergaster. Turkana Boy's brain was almost fully grown at the time of his death, but its volume (at 880 cc) was only about 130 cc greater than the maximum found in H. habilis, about 500 cc below the average of modern humans. The 130 cc increase from H. habilis becomes much less significant than what could be presumed when the larger body size of Turkana Boy and H. ergaster is considered. With all H. ergaster skulls considered, the brain volume of the species mostly varied between 600 and 910 cc, with some small examples only having a volume of 508–580 cc. Since their brain was smaller than that of modern humans, the skull of H. ergaster immediately narrowed behind the eye sockets (post-orbital constriction).
The brain case was long and low, and Turkana Boy's forehead was flat and receding, merging at an angle with the brow ridge above their eyes. A noticeable difference between Turkana Boy and the australopithecines and H. habilis would have been their nose, which would have been similar to that of modern humans in projecting forwards and having nostrils oriented downwards. This external nose may have also been an adaptation towards a warmer climate, since the noses of modern humans are usually cooler than their central bodies, condensing moisture that would otherwise have been exhaled and lost during periods of increased activity. The face of Turkana Boy would have been longer from top to bottom than that of modern humans, with the jaws projecting farther outwards (prognathism). Though the jaws and teeth were smaller than those of the average australopithecine and H. habilis, they were still significantly larger than those of modern humans. Since the jaw slanted sharply backwards, it is probable that they were chinless.
The overall structure of Turkana Boy's skull and face is also reflected in other H. ergaster skulls, which combine large and outwardly projecting faces with brow ridges, receding foreheads, large teeth and projecting nasal bones. Though Turkana Boy would have been no more than 12 years old when he died, their stature is more similar to that of a modern 15-year-old and the brain is comparable to that of a modern 1-year-old. By modern standards, H. ergaster would thus have been cognitively limited, though the invention of new tools prove that they were more intelligent than their predecessors.
Body mass and sexual dimorphism
H. ergaster possessed a significantly larger body mass in comparison to earlier hominins such as early Homo, Australopithecus and Paranthropus. Whereas australopithecines typically ranged in weight from 29–48 kg (64–106 lbs), H. ergaster typically ranged in weight from 52–63 kg (115–139 lbs). It is possible that the increased body size was the result of life in an open savannah environment, where increased size gives the ability to exploit broader diets in larger foraging areas, increases mobility and also gives the ability to hunt larger prey. The increased body mass also means that parents would have been able to carry their children to an older age and larger mass.
Though reduced sexual dimorphism has often been cited historically as one of the radical differences between H. ergaster and earlier Homo and australopithecines, it is unclear whether australopithecines were significantly more sexually diamorphic than H. ergaster or modern humans. Skeletal evidence suggests that sexes in H. ergaster differed no more in size than sexes in modern humans do, but a 2003 study by palaeoanthropologists Philip L. Reno, Richard S. Meindl, Melanie A. McCollum and C. Owen Lovejoy suggested that the same was also true for the significantly earlier Australopithecus afarensis. Sexual dimorphism is difficult to measure in extinct species since the sex of fossils is usually not determinable. Historically, scientists have typically measured differences between the extreme ends (in terms of size and morphology) of the fossil material attributed to a species and assumed that the resulting ratio applies to the mean difference between male and female individuals.
Growth and development
The dimensions of a 1.8 million years old adult female H. ergaster pelvis from Gona, Ethiopia suggests that H. ergaster would have been capable of birthing children with a maximum prenatal (pre-birth) brain size of 315 cc, about 30–50 % of adult brain size. This value falls intermediately between that of chimpanzees (~40 %) and modern humans (28%). Further conclusions about the growth and development in early Homo can be drawn from the Mojokerto child, a ~1.4–1.5 million year old ~1-year old Asian H. erectus, which had a brain at about 72–84% the size of an adult H. erectus brain, which suggests a brain growth trajectory more similar to that of other great apes than of modern humans. Both the Gona pelvis and the Mojokerto child suggest that the prenatal growth of H. ergaster was similar to that of modern humans but that the postnatal (post-birth) growth and development was intermediate between that of chimpanzees and modern humans. The faster development rate suggests that altriciality (an extended childhood and a long period of dependency on your parents) evolved at a later stage in human evolution, possibly in the last common ancestor of Neanderthals and modern humans. The faster development rate might also indicate that the expected lifespan of H. ergaster and H. erectus was lower than that of later and modern humans.
Culture
Diet and energetics
It is frequently assumed that the larger body and brain size of H. ergaster, compared to its ancestors, would have brought with it increased dietary and energy needs. In 2002, palaeoanthropologists Leslie C. Aiello and Jonathan C. K. Wells stated that the average resting metabolic requirements of H. ergaster would have been 39% higher than those of Australopithecus afarensis, 30% higher in males and 54% higher in females. However, the torso proportions of H. ergaster implies a relatively small gut, which means that energy needs might not necessarily have been higher in H. ergaster than in earlier hominins. This is because the earlier ape (and australopithecine) gut was large and energy-expensive since it needed to synthesize fat through fermenting plant matter, whereas H. ergaster likely ate significantly more animal fat than their predecessors. This would have allowed more energy to be diverted to brain growth, increasing brain size while maintaining the energy requirements of earlier species.
If they had increased energy requirements, H. ergaster would have needed to eat either vastly more food than australopithecines, or would have needed to eat food of superior quality. If they ate the same type of foods as the australopithecines, feeding time would then have had to be dramatically increased in proportion to the extra calories required, reducing the time H. ergaster could use for resting, socialising and travelling. Though this would have been possible, it is considered unlikely, especially since the jaws and teeth of H. ergaster are reduced in size compared to those of the australopithecines, suggesting a shift in diet away from fibrous and difficult-to-chew foods. Regardless of energy needs, the small gut of H. ergaster also suggests a more easily digested diet composed of food of higher quality.
It is likely that H. ergaster consumed meat in higher proportions than the earlier australopithecines. Meat was probably acquired through a combination of ambushes, active hunting and confrontational scavenging. H. ergaster must not only have possessed the ability of endurance running, but must also have been able to defend themselves and the carcasses of their prey from the variety of contemporary African predators. It is possible that a drop in African carnivoran species variety around 1.5 million years ago can be ascribed to competition with opportunistic and carnivorous hominins.
On its own, meat might not have been able to fully sustain H. ergaster. Modern humans can not sufficiently metabolize protein to meet more than 50% of their energy needs and modern humans who heavily rely on animal-based products in their diet mostly rely on fat to sustain the rest of their energy requirements. Multiple reasons make a fully meat-based diet in H. ergaster unlikely, the most prominent being that African ungulates (the primary prey available) are relatively low in fat and that high meat diets demand increased intake of water, which would have been difficult in an open and hot environment. Modern African hunter-gatherers who rely heavily on meat, such as the Hadza and San peoples, also use cultural means to recover the maximum amount of fat from the carcasses of their prey, a method that would not have been available to H. ergaster.
H. ergaster would thus likely have consumed large quantities of meat, vastly more than their ancestors, but would also have had to make use of a variety of other food sources, such as seeds, honey, nuts, invertebrates, nutritious tubers, bulbs and other underground plant storage organs. The relatively small chewing capacity of H. ergaster, in comparison to its larger-jawed ancestors, means that the meat and high quality plant food consumed would likely have required the use of tools to process before eating.
Social structure and dynamics
H. ergaster lived on the African savannah, which during the Pleistocene was home to a considerably more formidable community of carnivorans than the present savannah. Hominins could probably only have adapted to life on the savannah if effective anti-predator defense behaviours had already evolved. Defense against predators would likely have come through H. ergaster living in large groups, possessing stone (and presumably wooden) tools and effective counter-attack behaviour having been established. In modern primates that spend significant amounts of time on the savannah, such as chimpanzees and savannah baboons, individuals form large, multi-male, groups wherein multiple males can effectively work together to fend off and counter-attack predators, occasionally with the use of stones or sticks, and protect the rest of the group. It is possible that similar behaviour was exhibited in early Homo. Based on the male-bonded systems within bonobos and chimpanzees, and the tendency towards male bonding in modern foragers, groups of early Homo might have been male-bonded as well. Because of the scarcity of fossil material, group size in early Homo cannot be determined with any certainty. Groups were probably large, it is possible groups were above the upper range of known group sizes among chimpanzees and baboons ( 100 individuals or more). In 1993, palaeoanthropologists Leslie C. Aiello and R. I. M. Dunbar estimated that the group size of H. habilis and H. rudolfensis, based on neocortex size (as there is a known relationship between neocortex size and group size in modern non-human primates), would have ranged from about 70–85 individuals. With the additional factor of bipedalism, which is energetically cheaper than quadrupedalism, the maximum ecologically tolerable group size may have been even larger. Aiello's and Dunbar's group size estimate in regards to H. ergaster was 91–116 individuals.
Social and counter-attack behaviour of earlier Homo probably carried over into H. ergaster, where they are likely to have developed even further. H. ergaster was probably the first primate to move into the niche of social carnivore (i. e. hunter-gatherer). Such behaviour would probably have been the result of counter-attacks in the context of competition over nutritious food with other carnivores and would probably have evolved from something akin to the opportunistic hunting sometimes exhibited by chimpanzees. The switch to predation in groups might have triggered a cascade of evolutionary changes which changed the course of human evolution. Cooperative behaviours such as opportunistic hunting in groups, predator defense and confrontational scavenging would have been critical for survival which means that a fundamental transition in psychology gradually transpired. With the typical "competitive cooperation" behaviour exhibited by most primates no longer being favored through natural selection and social tendencies taking its place, hunting, and other activities, would have become true collaborative efforts. Because counter-attack behaviour is typically exhibited in males of modern primates, social hunting in archaic humans is believed to have been a primarily male activity. Females likely conducted other types of foraging, gathering food which did not require hunting (i.e. fruits, nuts, eggs etc.).
With hunting being a social activity, individuals probably shared the meat with one another, which would have strengthened the bonds both between the hunters themselves and between the hunters and the rest of the H. ergaster group. Females likely shared what they had foraged with the rest of the group as well. This development could have led to the development of male-female friendships into opportunistic monogamous pair bonds. Since sexual selection from females probably favored males that could hunt, the emerging social behaviour resulting from these new behaviours would have been carried over and amplified through the generations.
The only direct evidence of H. ergaster group composition comes from a series of sites outside of Ileret in Kenya, where 97 footprints made around 1.5 million years ago by a group of at least 20 individuals have been preserved. Based on the size of the footprints, one of the trackways appears to have been a group entirely composed of males, possibly a specialised task group, such as a border patrol or a hunting or foraging party. If this assessment is correct, this would further suggest a male-female division of responsibilities. In modern hunter-gatherer societies who target large prey items, male parties are typically dispatched to bring down these high-risk animals, and, due to the low success rate, female parties tend to focus on more predictable foods.
Technology
Tool production
Early H. ergaster inherited the Oldowan culture of tools from australopithecines and earlier Homo, though quickly learnt to strike much larger stone flakes than their predecessors and contemporaries. By 1.65 million years ago, H. ergaster had created the extensively flaked artefacts and early hand axes that mark the Acheulean culture, and by 1.6–1.4 million years ago, the new tool industry was widely established in East Africa. Acheulean tools differ from Oldowan tools in that the core forms of the tools were clearly deliberate. Whereas the shape of the core forms in Oldowan tools, which were probably used mostly as hammers to crack bones for marrow, appears to not have mattered much, the hand axes of the Acheulean culture demonstrate an intent to produce narrow and sharp objects, typically in teardrop, oval or triangular shapes. Once in place, the Acheulean industry remained unchanged throughout H. ergaster's existence and later times, with tools produced near its end about 250,000 years ago not being significantly different from tools produced 1.65 million years ago.
The oldest Acheulean assemblages also preserve core forms similar to those in Oldowan tools, but there are no known true intermediate forms between the two, suggesting that the appearance of Acheulean tools was an abrupt and sudden development. The most significant development that led to the Acheulean tools was likely early hominins learning the ability to strike large flakes, up to 30 cm (1 ft) or more in length, from larger boulders, from which they could manufacture new tools such as hand axes. Though "hand axe" implies that all hand axes were used for chopping and were hand-held, they came in a variety of different shapes and size and probably served several different functions. Carefully shaped and symmetric examples may have been hurled at prey akin to modern discuses, more casually made examples may simply have served as portable sources for sharp flakes and some could have been used for scraping or chopping wood. Additionally, hand axes are effective butchering tools and were possibly also used for dismembering carcasses of large animals.
There are preserved hand axes that are too unwieldy and large to be used for any apparent practical purpose. The use of these larger hand axes, and for some discovered collections of hundreds of hand axes without obvious signs of use, is speculative and conjectural. An idea that has been popular in the popular press, and frequently cited in academia, is that large and impressive hand axes might have been emblems used for attracting mates, with makers of large axes showing strength, coordination and determination, qualities that may have been regarded as attractive. Palaeoanthropologists April Nowell and Melanie Lee Chang noted in 2009 that though this theory is "both intriguing and emotionally appealing", there is little evidence for it and it is untestable. They considered it more probable that variations in hand axe morphology over the course of hundreds of thousands of years was the result of various different factors rather than a single, overarching factor in sexual selection.
Fire
As Homo migrated into open savannah environments, encounters with natural fires must have become more frequent and significant. It is possible that H. ergaster were the earliest humans to master the control of fire, which they may have used for cooking purposes. Cooking renders both meat and plant foods more digestible, which might have been important since the guts of H. ergaster were reduced in size compared to those of their ancestors. Though H. ergaster/H. erectus is frequently assumed to have been the earliest Homo to control fire, concrete evidence is somewhat lacking in the fossil record, perhaps partly due to the difficulty for actual evidence of fire usage to be preserved. Two of the earliest sites commonly claimed to preserve evidence of fire usage are FxJj20 at Koobi Fora and GnJi 1/6E near Lake Baringo, both in Kenya and both dated as up to 1.5 million years old. The evidence at FxJj20 consists of burned sediments and heat-altered stone tools, whereas GnJi 1/6E preserves large clasts of baked clay, associated with stone tools and faunal remains. Though it is difficult to exclude a natural origin for the fire residue evidenced, the sites remain strong candidates for early fire use.
Several sites, preserving more widely accepted evidence of fire usage, have been dated to 1 million years ago or younger, postdating the emergence and last generally accepted record of H. ergaster. These sites include cave sites, such as Wonderwerk and Swartkrans in South Africa, and open sites, such as Kalambo Falls in Zambia. The site Gesher Benot Ya’aqov in Israel, dated to about 700,000 years ago, preserves widely accepted evidence of fire usage through burnt materials and burnt flint microartefacts being preserved at numerous levels. From around 400,000 years ago and onwards, traces of fire become even more numerous in sites across Africa, Europe and Asia.
Language
The spinal cord of Turkana Boy would have been narrower than that of modern humans, which means that the nervous system of H. ergaster, and their respiratory muscles, may not have been developed enough to produce or control speech. In 2001, anthropologists Bruce Latimer and James Ohman concluded that Turkana Boy was afflicted by skeletal dysplasia and scoliosis, and thus would not have been representative of the rest of his species in this respect. In 2006, when anthropologist Marc Meyer and colleagues described a H. erectus s.l. specimen from Dmanisi, Georgia, dated to 1.78 million years old. The fossil preserves the oldest known Homo vertebrae and the spine found falls within the range of modern human spines, suggesting that the individual would have been capable of speech. Meyer and colleagues concluded that speech was probably possible within Homo very early on and that Turkana Boy probably suffered from some congenital defect, possibly spinal stenosis.
In 2013 and 2014, anthropologist Regula Schiess and colleagues concluded that there was no evidence of any congenital defects in Turkana Boy, and, in contrast to the 2001 and 2006 studies, considered the specimen to be representative of the species.
| Biology and health sciences | Homo | Biology |
276846 | https://en.wikipedia.org/wiki/Onchocerciasis | Onchocerciasis | Onchocerciasis, also known as river blindness, is a disease caused by infection with the parasitic worm Onchocerca volvulus. Symptoms include severe itching, bumps under the skin, and blindness. It is the second-most common cause of blindness due to infection, after trachoma.
The parasitic worm is spread by the bites of a black fly of the Simulium genus. Usually, many bites are required before infection occurs. These flies live near rivers, hence the common name of the disease. Once inside a person, the worms create larvae that make their way out to the skin, where they can infect the next black fly that bites the person. There are a number of ways to make the diagnosis, including placing a biopsy of the skin in normal saline and watching for the larva to come out, looking in the eye for larvae, and looking within the bumps under the skin for adult worms.
A vaccine against the disease does not exist. Prevention is by avoiding being bitten by flies. This may include the use of insect repellent and proper clothing. Other efforts include those to decrease the fly population by spraying insecticides. Efforts to eradicate the disease by treating entire groups of people twice a year are ongoing in a number of areas of the world. Treatment of those infected is with the medication ivermectin every six to twelve months. This treatment kills the larvae but not the adult worms. The antibiotic doxycycline weakens the worms by killing an associated bacterium, Wolbachia, and is recommended by some as well. The lumps under the skin may also be removed by surgery.
According to the Centers for Disease Control and Prevention, as of 2017, about 20.9 million people were infected with onchocerciasis, and an estimated 1.15 million have some vision loss from the infection. Most infections occur in sub-Saharan Africa, although cases have also been reported in Yemen and isolated areas of Central and South America. In 1915, the physician Rodolfo Robles first linked the worm to eye disease. It is listed by the World Health Organization (WHO) as a neglected tropical disease. In 2013 Colombia became the first country to eradicate the disease.
Cause
Onchocerciasis is a parasitic infection caused by the roundworm species Onchocerca volvulus. The larvae of O. volvulus enter a human host when an infected female adult fly from the genus Simulium bites them. After that, it can take up to three months for the worms to mature under the skin of their host. The worms mainly get nutrients for growth in humans from blood, but they have also been seen to rely on other bodily fluids such as cerebrospinal fluid, and urine. It is common to see nodules formed in the skin where the adult worms reside and mate. However, these worms will often travel throughout the body using blood vessels in connective tissues and even settle behind the cornea.
Life cycle
The life of the parasite can be traced through the black fly and the human hosts in the following steps:
A Simulium female black fly takes a blood meal on an infected human host and ingests microfilaria.
The microfilaria enter the black fly's gut and thoracic flight muscles, progressing into the first larval stage (J1.).
The larvae mature into the second larval stage (J2.) and move to the proboscis and into the saliva in its third larval stage (J3.). Maturation takes about seven days.
The black fly takes another blood meal, passing the larvae into the next human host's blood.
The larvae migrate to the subcutaneous tissue and undergo two more molts. They form nodules as they mature into adult worms over six to 12 months.
After maturing, adult male worms mate with female worms in the subcutaneous tissue to produce between 700 and 1,500 microfilariae daily.
The microfilaria migrate to the skin during the day, and the black flies only feed during the day, so the parasite is in a prime position for the female fly to ingest it. Black flies take blood meals to ingest these microfilaria to restart the cycle.
Signs and symptoms
The larvae can move through the body without triggering a response from the host's immune system, so some people who are infected with the parasite experience no symptoms; the Global Burden of Disease Study estimated that in 2017 there were at least 20.9 million people infected worldwide, of which 14.6 million had skin disease symptoms and 1.15 million experienced symptoms that impacted vision. After a blackfly bite, it can take 12–18 months for the larvae to develop into mature adult worms that will produce their larvae, which is what leads to the development of symptoms. Almost all the clinical manifestations of onchocerciasis are due to localized host inflammatory responses to dead or dying microfilariae (larvae). The signs and symptoms of onchocerciasis are usually divided into two categories, skin and eye symptoms.
Skin symptoms will develop years before any vision problems. These symptoms include:
Intense itching
Swelling
Inflammation
Depigmentation
Hyperpigmentation
Rash
Nodules under the skin
Skin atrophy
Hanging groin (folds of inelastic atrophic skin in the groin associated with enlarged lymph nodes)
Eye symptoms include:
Vision impairment, low vision, or permanent blindness.
Clouding of the cornea
Light sensitivity
Lesions on eyes
Glaucoma
Eye pain
Eye redness
Eye symptoms provide the common name associated with onchocerciasis, river blindness, and may involve any part of the eye from conjunctiva and cornea to uvea and posterior segment, including the retina and optic nerve. The microfilariae migrate to the surface of the cornea. Punctate keratitis occurs in the infected area. This clears up as the inflammation subsides. However, if the infection is chronic, sclerosing keratitis can occur, making the affected area become opaque. Over time, the entire cornea may become opaque, thus leading to blindness. Some evidence suggests an immune response to bacteria present in the worms causes the effect on the cornea.
Mazzotti reaction
The Mazzotti reaction, first described in 1948, is a symptom complex seen in patients after undergoing treatment of onchocerciasis with the medication diethylcarbamazine (DEC). Mazzotti reactions can be life-threatening, and are characterized by fever, urticaria, swollen and tender lymph nodes, tachycardia, hypotension, arthralgias, oedema, and abdominal pain that occur within seven days of treatment of microfilariasis.
The phenomenon is so common when DEC is used that this drug is the basis of a skin patch test used to confirm that diagnosis. The drug patch is placed on the skin, and if the patient is infected with O. volvulus microfilaria, localized pruritus and urticaria are seen at the application site.
Nodding disease
This is an unusual form of epidemic epilepsy associated with onchocerciasis although a definitive link has not been established. This syndrome was first described in Tanzania by Louise Jilek-Aall, a Norwegian psychiatric doctor in Tanzanian practice, during the 1960s. It occurs most commonly in Uganda and South Sudan. It manifests itself in previously healthy 5–15-year-old children, is often triggered by eating or low temperatures, and is accompanied by cognitive impairment. Seizures occur frequently and may be difficult to control. The electroencephalogram is abnormal but cerebrospinal fluid (CSF) and magnetic resonance imaging (MRI) are normal or show non-specific changes. If there are abnormalities on the MRI, they are usually present in the hippocampus.
Diagnosis
When a clinical diagnosis of onchocerciasis is obtained, doctors take small snips of skin containing 3–5 mg of skin tissue. The skin samples taken are only from the upper dermis. These samples will then be soaked in saline and examined underneath a microscope to check for the presence of microfilaria. If microfilariae are not detected in the samples, the Mazzotti test is used. In this test, 6 mg of diethylcarbamazine is administered to the affected area. If the patient experiences intense inflammation or itching in the affected area Microfilaria is present. Slit-lamp eye exams are used to identify signs of the parasites in and around the eyes of patients whose eyes are affected. Antibody tests when available can aid in the diagnosis of Onchocerciasis.
Classification
Onchocerciasis causes different kinds of skin changes, which vary in different geographic regions; it may be divided into the following phases or types:
Erisipela de la costa
An acute phase, it is characterized by swelling of the face, with erythema and itching. This skin change, erisípela de la costa, of acute onchocerciasis is most commonly seen among victims in Central and South America.
Mal morando
This cutaneous condition is characterized by inflammation accompanied by hyperpigmentation.
Sowda
A cutaneous condition, it is a localized type of onchocerciasis.
Additionally, the various skin changes associated with onchocerciasis may be described as follows:
Leopard skin
The spotted depigmentation of the skin that may occur with onchocerciasis
Elephant skin
The thickening of human skin that may be associated with onchocerciasis
Lizard skin
The thickened, wrinkled skin changes that may result with onchocerciasis
Prevention
Various control programs aim to stop onchocerciasis from being a public health problem. The Onchocerciasis Control Programme (OCP) launched in 1974, and at its peak, covered 30 million people in the following countries: Benin, Burkina Faso, Côte d'Ivoire, Ghana, Togo, Mali, and Niger. The OCP utilized the following initiatives: the use of larvicide spraying into fast-flowing rivers to control black fly populations, and from 1988 onwards, the use of ivermectin to treat infected people as a core treatment therapy. Alongside the OCP, a joint effort of the World Health Organization, the World Bank, the United Nations Development Programme, and the UN Food and Agriculture Organization, was considered to be a success in controlling onchocerciasis, and in 2002 shifted from control of onchocerciasis to elimination. According to the World Health Organization, four countries have eradicated onchocerciasis: Colombia (2013), Ecuador (2014), Mexico (2015), and Guatemala (2016). Continued monitoring ensures onchocerciasis cannot reinvade the area through the OCP. Other effective prevention efforts include personal protection from black fly bites. Recommended protection measures from the CDC include using insect repellents and wearing long sleeves and pants to eliminate exposed skin. Using insect repellent that contains N,N-Diethyl-meta-toluamide (DEET) as well as clothing treated with permethrin.
Elimination
In 1995, the African Programme for Onchocerciasis Control (APOC) was initiated to eliminate onchocerciasis in African countries where the disease was endemic. The initiative relied primarily on the use of the antiparasitic drug ivermectin. The initiative was to set up community-directed treatment with ivermectin for those at risk of infection. Overall transmission has declined. The APOC ended in 2015 and aspects of its work has been taken over by the WHO Expanded Special Programme for the Elimination of Neglected Tropical Diseases (ESPEN). As in the Americas, the objective of ESPEN has been to work with Government Health Ministries and partner non-governmental organizations, to eliminate the transmission of onchocerciasis. This requires consistent annual treatment of 80% of the population in endemic areas for at least 10–12 years, the life span of the adult worm. No African country has verified the elimination of onchocerciasis. Treatment has stopped in some areas (e.g. Nigeria), following epidemiological and entomological assessments that demonstrated no ongoing transmission could be detected.
The Onchocerciasis Elimination Programme for the Americas (OEPA) was launched in 1992. On July 29, 2013, the Pan American Health Organization (PAHO) announced that after 16 years of efforts, Colombia had become the first country in the world to eliminate onchocerciasis. Countries that received verification of elimination were Colombia in 2013, Ecuador in 2015, and Guatemala in 2016. The key factor in elimination was mass administration of ivermectin. The OEPA projection was that the disease would be eliminated from all remaining countries in the Americas by 2012. In September 2015, the OEPA announced that onchocerciasis only remained in a remote region on the border of Brazil and Venezuela. The area is home to the Yanomami indigenous people.
No vaccine to prevent onchocerciasis infection in humans is available. This is due to two potential target product profiles (TPPs) that must be considered when developing a vaccine for onchocerciasis, making vaccine development difficult. The project to develop a vaccine for the disease has two goals. The first priority is a preventive vaccine for children five years or younger, as this population does not receive ivermectin. The second is a therapeutic vaccine to target adult worms, microfilariae, and the causative agents of pathology and transmission, or both, for children and adults with O. volvulus infection.
Treatment
In mass drug administration (MDA) programmes, the treatment for onchocerciasis is ivermectin (trade name: Mectizan). Ivermectin is administered four times a year and will be continually administered for 10–14 years due to the lifespan of the adult worm. Intense skin itching is eventually relieved, and the progression towards blindness is halted. The drug reduces larval release from the adult worm but does not kill it. However the drug does not prevent transmission of Onchocerciasis. It however reduces morbidity and has shown promising results to eliminate in some endemic areas of Africa
Ivermectin treatment is particularly effective because it only needs to be taken once or twice a year, needs no refrigeration, and has a wide margin of safety, with the result that it has been widely given by minimally trained community health workers. Patients taking the drug for the treatment of onchocerciasis may have adverse effects within 1–2 days after the drug is administered. Symptoms include urticaria, pruritus, fever, dermatitis, myalgia, swelling of the face and limbs, or postural hypotension.
Antibiotics
For the treatment of individuals, doxycycline is used to kill the Wolbachia bacteria that live in adult worms. This adjunct therapy has been shown to significantly lower microfilarial loads in the host and may kill the adult worms, due to the symbiotic relationship between Wolbachia and the worm. In four separate trials over ten years with various dosing regimens of doxycycline for individualized treatment, doxycycline was found to be effective in sterilizing the female worms and reducing their numbers over four to six weeks. Research on other antibiotics, such as rifampicin, has shown it to be effective in animal models at reducing Wolbachia both as an alternative and as an adjunct to doxycycline. However, doxycycline treatment requires daily dosing for at least four to six weeks, making it more difficult to administer in the affected areas.
Ivermectin
Ivermectin kills the parasite by interfering with the nervous system and muscle function, in particular, by enhancing inhibitory neurotransmission. The drug binds to and activates glutamate-gated chloride channels. These channels, present in neurons and myocytes, are not invertebrate-specific, but are protected in vertebrates from the action of ivermectin by the blood–brain barrier. Ivermectin is thought to irreversibly activate these channel receptors in the worm, eventually causing an inhibitory postsynaptic potential. The chance of a future action potential occurring in synapses between neurons decreases and the nematodes experience flaccid paralysis followed by death.
Ivermectin is directly effective against the larval stage microfilariae of O. volvulus; they are paralyzed and can be killed by eosinophils and macrophages. It does not kill adult females (macrofilariae) but does cause them to cease releasing microfilariae, perhaps by paralyzing the reproductive tract. Ivermectin is very effective in reducing microfilarial load and reducing number of punctate opacities in individuals with onchocerciasis.
Moxidectin
After two decades of research, moxidectin was approved by the U.S. Food and Drug Administration in 2018 for use in ages 12 and older. Ongoing studies are looking to identify doses that will be safe for children ages 4–11. The oral dosage for moxidectin in adults and children 12 and up is 8 mg in a single dose. Moxidectin has been found to more strongly suppress the O. volvulus microfilariae for longer than ivermectin treatments, with peak clearance of microfilariae in the skin at one month after treatment. After six months post-treatment, many individuals treated with moxidectin have no detectable microfilariae in their skin.
Epidemiology
About 21 million people were infected with this parasite in 2017; about 1.2 million of those had vision loss. As of 2017, about 99% of onchocerciasis cases occurred in Africa. Onchocerciasis is currently relatively common in 31 African countries, Yemen, and isolated regions of South America. Over 85 million people live in endemic areas, and half of these reside in Nigeria. Another 120 million people are at risk for contracting the disease. The Onchocerca volvulus main habitat is fast-flowing rivers, Onchocerciasis is more commonly found along the large rivers in northern and central regions of Africa, with cases decreasing with distance from the rivers. Multiple exposures to Simulium blackflies raise the number of adult worms and microfilariae that are present in the host. Risk of contracting Onchocerciasis for casual travelers is low, since it often takes several exposures, while travelers that stay for longer visits such as missionaries or long-term volunteers have a greater risk of contracting Onchocerciasis. Onchocerciasis was eliminated in the northern focus in Chiapas, Mexico, and the focus in Oaxaca, Mexico, where Onchocerca volvulus existed, was determined, after several years of treatment with ivermectin, as free of the transmission of the parasite. In April 2013, Colombia became the first country to achieve elimination of Onchocerciasis, verified by the World Health Organization. In the following three years, Ecuador, Guatemala, and Mexico eliminated Onchocerciasis with ivermectin.
Cities in Nigeria, Cameroon, Ethiopia, Uganda, and the Congo have had the largest amounts of infected individuals.
The efforts of CDTI (Community-Directed Treatment with Ivermectin) were conducted to study Onchocerciasis associations with epilepsy. The results do not go unnoticed as they decreased the number of microfilariae (larvae) loads. This was able to decrease the number of blind people due to onchocerciasis dramatically. However, another issue that arises is the fact that onchocerciasis can cause epilepsy, most likely because the level of microfilariae load required to develop epilepsy is much lower than to develop blindness.
According to a 2002 WHO report, onchocerciasis has not caused a single death, but its global burden is 987,000 disability adjusted life years (DALYs). The severe itchiness alone accounts for 60% of the DALYs. Infection reduces the host's immunity and resistance to other diseases, which results in an estimated reduction in life expectancy of 13 years. In 2017, the Global Burden of Disease study said that an estimated 220 million people needed preventive chemotherapy for onchocerciasis. Of those infected, 14.6 million had skin disease, and 1.15 million experienced vision loss.
Onchocerciasis is the second leading cause of blindness from infectious causes. Main disease symptoms, such as blindness and itching, contribute to disease burden by limiting the infected individuals' ability to live and work. Individuals most at risk live or work in areas where Simulium blackflies are most common, mostly near rivers and streams. Rural agricultural areas in sub-Saharan Africa see the most disease burden by blackfly bites. Onchocerciasis common to tropical environments, like that of sub-Saharan Africa, where more than 99% percent of infected individuals occupy the 31 countries. Onchocerciasis can be linked to impoverished remote areas, as residents who experience symptoms can no longer tend to land or navigate the area. Areas with high infection rates may experience up to one-third of residents affected by onchocerciasis symptoms. The age group most impacted by the disease are individuals age 61+ years.
History
Onchocerca originated in Africa and was exported to the Americas by the slave trade, as part of the Columbian exchange that introduced other old-world diseases such as yellow fever into the New World. The findings of a phylogenetic study in the mid-90s are consistent with an introduction to the New World in this manner. DNA sequences of savannah and rainforest strains in Africa differ, while American strains are identical to savannah strains in western Africa. The microfilarial parasite that causes the disease was first identified in 1874 by an Irish naval surgeon, John O'Neill, who sought to identify the cause of a common skin disease along the west coast of Africa, known as "craw-craw". Rudolf Leuckart, a German zoologist, later examined specimens of the same filarial worm sent from Africa by a German missionary doctor in 1890 and named the organism Filaria volvulus.
Rodolfo Robles and Rafael Pacheco in Guatemala first mentioned the ocular form of the disease in the Americas in about 1915. They described a tropical worm infection with adult Onchocerca that included inflammation of the skin, especially the face ('erisipela de la costa'), and eyes. The disease, commonly called the "filarial blinding disease", and later referred to as "Robles disease", was common among coffee plantation workers. Manifestations included subcutaneous nodules, anterior eye lesions, and dermatitis. Robles sent specimens to Émile Brumpt, a French parasitologist, who named it O. caecutiens in 1919, indicating the parasite caused blindness (Latin "caecus" meaning blind). The disease was also reported as being common in Mexico. By the early 1920s, it was generally agreed that the filaria in Africa and Central America were morphologically indistinguishable and the same as that described by O'Neill 50 years earlier.
Robles hypothesized that the vector of the disease was the day-biting black fly, Simulium. Scottish physician Donald Blacklock of the Liverpool School of Tropical Medicine confirmed this mode of transmission in studies in Sierra Leone. Blacklock's experiments included the re-infection of Simulium flies exposed to portions of the skin of infected subjects on which nodules were present, which led to the elucidation of the life cycle of the Onchocerca parasite. Blacklock and others could find no evidence of eye disease in Africa. Jean Hissette, a Belgian ophthalmologist, discovered in 1930 that the organism was the cause of a "river blindness" in the Belgian Congo. Some of the patients reported seeing tangled threads or worms in their vision, which were microfilariae moving freely in the aqueous humor of the anterior chamber of the eye. Blacklock and Strong had thought the African worm did not affect the eyes, but Hissette reported that 50% of patients with onchocerciasis near the Sankuru River in the Belgian Congo had eye disease and 20% were blind. Hisette Isolated the microfilariae from an enucleated eye and described the typical chorioretinal scarring, later called the "Hissette-Ridley fundus" after another ophthalmologist, Harold Ridley, who also made extensive observations on onchocerciasis patients in northwest Ghana, publishing his findings in 1945. Ridley first postulated that the disease was brought by the slave trade. The international scientific community was initially skeptical of Hisette's findings, but they were confirmed by the Harvard African Expedition of 1934, led by Richard P. Strong, an American tropical medicine physician.
Society and culture
Since 1987, ivermectin has been provided free of charge for use in humans by Merck through the Mectizan donation program (MDP). The MDP works together with ministries of health and nongovernmental development organisations, such as the World Health Organization, to provide free ivermectin to those who need it in endemic areas. Due to the joint efforts of NGOs and WHO, onchocerciasis is no longer an obstacle in socio-economic development.
In 2015 William C. Campbell and Satoshi Ōmura were co-awarded half of that year's Nobel Prize in Physiology or Medicine for the discovery of the avermectin family of compounds, the forerunner of ivermectin. The latter has come to decrease the occurrence of lymphatic filariasis and onchocerciasis.
Uganda's government, working with the Carter Center river blindness program since 1996, switched strategies for the distribution of Mectizan. The male-dominated volunteer distribution system had failed to take advantage of traditional kinship structures and roles. In 2014, the program switched from village health teams to community distributors, primarily selecting women to assure that everyone in the circle of their family and friends received river blindness information and Mectizan.
In 2021, Nigeria had the greatest prevalence of onchocerciasis infections globally and attributed the infection to 30.2% of blindness cases in the country. A study from western Nigeria found that residents believed that the parasitic effects of the disease were necessary to stimulate fertility and that the disease was thought to be carried by all residents.
Research
Animal models for the disease are somewhat limited, as the parasite only lives in primates, but there are close parallels. Litomosoides sigmodontis , which will naturally infect cotton rats, has been found to fully develop in BALB/c mice. Onchocerca ochengi, the closest relative of O. volvulus, lives in intradermal cavities in cattle, and is also spread by black flies. Both systems are useful, but not exact, animal models.
A study of 2501 people in Ghana showed the prevalence rate doubled between 2000 and 2005 despite treatment, suggesting the parasite is developing resistance to the drug. A clinical trial of another anti-parasitic agent, moxidectin (manufactured by Wyeth), began on July 1, 2009 (NCT00790998).
A Cochrane review compared outcomes of people treated with ivermectin alone versus doxycycline plus ivermectin. While there were no differences in most vision-related outcomes between the two treatments, there was low-quality evidence suggesting treatment with doxycycline plus ivermectin showed improvement in iridocyclitis and punctate keratitis, over those treated with ivermectin alone.
In 2017, WHO set up the Onchocerciasis Technical Advisory Subgroup (OTS) to further research and establish areas that require drug administration. The OTS also identifies co-endemic areas with lymphatic filariasis to properly treat Onchocerciasis and lymphatic filariasis.
WHO prioritizes research to eliminate onchocerciasis. Research approaches include: improving outreach efforts to marginalized populations, expanding the mapping of endemic areas of onchocerciasis, improving and standardizing information on mass drug administration, develop diagnostic approaches, surveillance strategies, and therapeutic approaches.
| Biology and health sciences | Helminthic diseases and infestations | Health |
276881 | https://en.wikipedia.org/wiki/Extreme%20value%20theorem | Extreme value theorem | In calculus, the extreme value theorem states that if a real-valued function is continuous on the closed and bounded interval , then must attain a maximum and a minimum, each at least once.
That is, there exist numbers and in such that:
The extreme value theorem is more specific than the related boundedness theorem, which states merely that a continuous function on the closed interval is bounded on that interval; that is, there exist real numbers and such that:
This does not say that and are necessarily the maximum and minimum values of on the interval which is what the extreme value theorem stipulates must also be the case.
The extreme value theorem is used to prove Rolle's theorem. In a formulation due to Karl Weierstrass, this theorem states that a continuous function from a non-empty compact space to a subset of the real numbers attains a maximum and a minimum.
History
The extreme value theorem was originally proven by Bernard Bolzano in the 1830s in a work Function Theory but the work remained unpublished until 1930. Bolzano's proof consisted of showing that a continuous function on a closed interval was bounded, and then showing that the function attained a maximum and a minimum value. Both proofs involved what is known today as the Bolzano–Weierstrass theorem.
Functions to which the theorem does not apply
The following examples show why the function domain must be closed and bounded in order for the theorem to apply. Each fails to attain a maximum on the given interval.
defined over is not bounded from above.
defined over is bounded but does not attain its least upper bound .
defined over is not bounded from above.
defined over is bounded but never attains its least upper bound .
Defining in the last two examples shows that both theorems require continuity on .
Generalization to metric and topological spaces
When moving from the real line to metric spaces and general topological spaces, the appropriate generalization of a closed bounded interval is a compact set. A set is said to be compact if it has the following property: from every collection of open sets such that , a finite subcollection can be chosen such that . This is usually stated in short as "every open cover of has a finite subcover". The Heine–Borel theorem asserts that a subset of the real line is compact if and only if it is both closed and bounded. Correspondingly, a metric space has the Heine–Borel property if every closed and bounded set is also compact.
The concept of a continuous function can likewise be generalized. Given topological spaces , a function is said to be continuous if for every open set , is also open. Given these definitions, continuous functions can be shown to preserve compactness:
In particular, if , then this theorem implies that is closed and bounded for any compact set , which in turn implies that attains its supremum and infimum on any (nonempty) compact set . Thus, we have the following generalization of the extreme value theorem:
Slightly more generally, this is also true for an upper semicontinuous function. (see compact space#Functions and compact spaces).
Proving the theorems
We look at the proof for the upper bound and the maximum of . By applying these results to the function , the existence of the lower bound and the result for the minimum of follows. Also note that everything in the proof is done within the context of the real numbers.
We first prove the boundedness theorem, which is a step in the proof of the extreme value theorem. The basic steps involved in the proof of the extreme value theorem are:
Prove the boundedness theorem.
Find a sequence so that its image converges to the supremum of .
Show that there exists a subsequence that converges to a point in the domain.
Use continuity to show that the image of the subsequence converges to the supremum.
Proof of the boundedness theorem
Proofs of the extreme value theorem
Proof using the hyperreals
Proof from first principles
Statement If is continuous on then it attains its supremum on
Extension to semi-continuous functions
If the continuity of the function f is weakened to semi-continuity,
then the corresponding half of the boundedness theorem and the extreme value theorem hold and the values –∞ or +∞, respectively, from the extended real number line can be allowed as possible values.
A function is said to be upper semi-continuous if
Applying this result to −f proves a similar result for the infimums of lower semicontinuous functions.
A function is said to be lower semi-continuous if
A real-valued function is upper as well as lower semi-continuous, if and only if it is continuous in the usual sense. Hence these two theorems imply the boundedness theorem and the extreme value theorem.
| Mathematics | Real analysis | null |
276918 | https://en.wikipedia.org/wiki/Division%20sign | Division sign | The division sign () is a mathematical symbol consisting of a short horizontal line with a dot above and another dot below, used in Anglophone countries to indicate the operation of division. This usage, though widespread in some countries, is not universal and the symbol has a different meaning in other countries. Its use to denote division is not recommended in the ISO 80000-2 standard for mathematical notation.
In mathematics
The obelus, a historical glyph consisting of a horizontal line with (or without) one or more dots, was first used as a symbol for division in 1659, in the algebra book by Johann Rahn, although previous writers had used the same symbol for subtraction. Some near-contemporaries believed that John Pell, who edited the book, may have been responsible for this use of the symbol. Other symbols for division include the slash or solidus , the colon , and the fraction bar (the horizontal bar in a vertical fraction). The ISO 80000-2 standard for mathematical notation recommends only the solidus or "fraction bar" for division, or the "colon" for ratios; it says that the sign "should not be used" for division.
In Italy, Poland and Russia, the sign was sometimes used to denote a range of values, and in Scandinavian countries it was, and sometimes still is, used as a negation sign: the Unicode Consortium has allocated a separate code point, for this usage uniquely; the exact form of the symbol displayed is typeface (font) dependent.
In computer systems
Encoding
The symbol was assigned to code point 0xF7 in ISO 8859-1, as the "division sign". This encoding was transferred to Unicode as U+00F7. In HTML, it can be encoded as or (at HTML level 3.2), or as .
Unicode provides various division symbols:
| Mathematics | Basics | null |
277102 | https://en.wikipedia.org/wiki/Syndrome | Syndrome | A syndrome is a set of medical signs and symptoms which are correlated with each other and often associated with a particular disease or disorder. The word derives from the Greek σύνδρομον, meaning "concurrence". When a syndrome is paired with a definite cause this becomes a disease. In some instances, a syndrome is so closely linked with a pathogenesis or cause that the words syndrome, disease, and disorder end up being used interchangeably for them. This substitution of terminology often confuses the reality and meaning of medical diagnoses. This is especially true of inherited syndromes. About one third of all phenotypes that are listed in OMIM are described as dysmorphic, which usually refers to the facial gestalt. For example, Down syndrome, Wolf–Hirschhorn syndrome, and Andersen–Tawil syndrome are disorders with known pathogeneses, so each is more than just a set of signs and symptoms, despite the syndrome nomenclature. In other instances, a syndrome is not specific to only one disease. For example, toxic shock syndrome can be caused by various toxins; another medical syndrome named as premotor syndrome can be caused by various brain lesions; and premenstrual syndrome is not a disease but simply a set of symptoms.
If an underlying genetic cause is suspected but not known, a condition may be referred to as a genetic association (often just "association" in context). By definition, an association indicates that the collection of signs and symptoms occurs in combination more frequently than would be likely by chance alone.
Syndromes are often named after the physician or group of physicians that discovered them or initially described the full clinical picture. Such eponymous syndrome names are examples of medical eponyms. Recently, there has been a shift towards naming conditions descriptively (by symptoms or underlying cause) rather than eponymously, but the eponymous syndrome names often persist in common usage.
The defining of syndromes has sometimes been termed syndromology, but it is usually not a separate discipline from nosology and differential diagnosis generally, which inherently involve pattern recognition (both sentient and automated) and differentiation among overlapping sets of signs and symptoms. Teratology (dysmorphology) by its nature involves the defining of congenital syndromes that may include birth defects (pathoanatomy), dysmetabolism (pathophysiology), and neurodevelopmental disorders.
Subsyndromal
When there are a number of symptoms suggesting a particular disease or condition but does not meet the defined criteria used to make a diagnosis of that disease or condition. This can be a bit subjective because it is ultimately up to the clinician to make the diagnosis. This could be because it has not advanced to the level or passed a threshold or just similar symptoms cause by other issues. Subclinical is synonymous since one of its definitions is "where some criteria are met but not enough to achieve clinical status"; but subclinical is not always interchangeable since it can also mean "not detectable or producing effects that are not detectable by the usual clinical tests"; i.e., asymptomatic.
Usage
General medicine
In medicine, a broad definition of syndrome is used, which describes a collection of symptoms and findings without necessarily tying them to a single identifiable pathogenesis. Examples of infectious syndromes include encephalitis and hepatitis, which can both have several different infectious causes. The more specific definition employed in medical genetics describes a subset of all medical syndromes.
History
Early texts by physicians noted the symptoms of various maladies and introduced diagnoses based upon those symptoms. For example, Avicenna's The Canon of Medicine (1025) describes diagnosing pleurisy by its symptoms, including chronic fever, cough, shooting pains, and labored breathing. The 17th century doctor Thomas Sydenham likewise approached diagnoses based upon collections of symptoms.
Psychiatry and psychopathology
Psychiatric syndromes often called psychopathological syndromes (psychopathology refers both to psychic dysfunctions occurring in mental disorders, and the study of the origin, diagnosis, development, and treatment of mental disorders).
In Russia those psychopathological syndromes are used in modern clinical practice and described in psychiatric literature in the details: asthenic syndrome, obsessive syndrome, emotional syndromes (for example, manic syndrome, depressive syndrome), Cotard's syndrome, catatonic syndrome, hebephrenic syndrome, delusional and hallucinatory syndromes (for example, paranoid syndrome, paranoid-hallucinatory syndrome, Kandinsky-Clérambault's syndrome also known as syndrome of psychic automatism, hallucinosis), paraphrenic syndrome, psychopathic syndromes (includes all personality disorders), clouding of consciousness syndromes (for example, twilight clouding of consciousness, amential syndrome also known as amentia, delirious syndrome, stunned consciousness syndrome, oneiroid syndrome), hysteric syndrome, neurotic syndrome, Korsakoff's syndrome, hypochondriacal syndrome, paranoiac syndrome, senestopathic syndrome, encephalopathic syndrome.
Some examples of psychopathological syndromes used in modern Germany are psychoorganic syndrome, depressive syndrome, paranoid-hallucinatory syndrome, obsessive-compulsive syndrome, autonomic syndrome, hostility syndrome, manic syndrome, apathy syndrome.
Münchausen syndrome, Ganser syndrome, neuroleptic-induced deficit syndrome, olfactory reference syndrome are also well-known.
History
The most important psychopathological syndromes were classified into three groups ranked in order of severity by German psychiatrist Emil Kraepelin (1856—1926). The first group, which includes the mild disorders, consists of five syndromes: emotional, paranoid, hysterical, delirious, and impulsive. The second, intermediate, group includes two syndromes: schizophrenic syndrome and speech-hallucinatory syndrome. The third includes the most severe disorders, and consists of three syndromes: epileptic, oligophrenic and dementia. In Kraepelin's era, epilepsy was viewed as a mental illness; Karl Jaspers also considered "genuine epilepsy" a "psychosis", and described "the three major psychoses" as schizophrenia, epilepsy, and manic-depressive illness.
Medical genetics
In the field of medical genetics, the term "syndrome" is traditionally only used when the underlying genetic cause is known. Thus, trisomy 21 is commonly known as Down syndrome.
Until 2005, CHARGE syndrome was most frequently referred to as "CHARGE association". When the major causative gene (CHD7) for the condition was discovered, the name was changed. The consensus underlying cause of VACTERL association has not been determined, and thus it is not commonly referred to as a "syndrome".
Other fields
In biology, "syndrome" is used in a more general sense to describe characteristic sets of features in various contexts. Examples include behavioral syndromes, as well as pollination syndromes and seed dispersal syndromes.
In orbital mechanics and astronomy, Kessler syndrome refers to the effect where the density of objects in low Earth orbit (LEO) is high enough that collisions between objects could cause a cascade in which each collision generates space debris that increases the likelihood of further collisions.
In quantum error correction theory syndromes correspond to errors in code words which are determined with syndrome measurements, which only collapse the state on an error state, so that the error can be corrected without affecting the quantum information stored in the code words.
Naming
There is no set common convention for the naming of newly identified syndromes. In the past, syndromes were often named after the physician or scientist who identified and described the condition in an initial publication. These are referred to as "eponymous syndromes". In some cases, diseases are named after the patient who initially presents with symptoms, or their home town (Stockholm syndrome). There have been isolated cases of patients being eager to have their syndromes named after them, while their physicians are hesitant. When a syndrome is named after a person, there is some difference of opinion as to whether it should take the possessive form or not (e.g. Down syndrome vs. Down's syndrome). North American usage has tended to favor the non-possessive form, while European references often use the possessive. A 2009 study demonstrated a trend away from the possessive form in Europe in medical literature from 1970 through 2008.
Underlying cause
Even in syndromes with no known etiology, the presence of the associated symptoms with a statistically improbable correlation normally leads the researchers to hypothesize that there exists an unknown underlying cause for all the described symptoms.
| Biology and health sciences | Miscellaneous | null |
277184 | https://en.wikipedia.org/wiki/Mathematical%20notation | Mathematical notation | Mathematical notation consists of using symbols for representing operations, unspecified numbers, relations, and any other mathematical objects and assembling them into expressions and formulas. Mathematical notation is widely used in mathematics, science, and engineering for representing complex concepts and properties in a concise, unambiguous, and accurate way.
For example, the physicist Albert Einstein's formula is the quantitative representation in mathematical notation of mass–energy equivalence.
Mathematical notation was first introduced by François Viète at the end of the 16th century and largely expanded during the 17th and 18th centuries by René Descartes, Isaac Newton, Gottfried Wilhelm Leibniz, and overall Leonhard Euler.
Symbols
The use of many symbols is the basis of mathematical notation. They play a similar role as words in natural languages. They may play different roles in mathematical notation similarly as verbs, adjective and nouns play different roles in a sentence.
Letters as symbols
Letters are typically used for naming—in mathematical jargon, one says representing—mathematical objects. The Latin and Greek alphabets are used extensively, but a few letters of other alphabets are also used sporadically, such as the Hebrew , Cyrillic , and Hiragana . Uppercase and lowercase letters are considered as different symbols. For Latin alphabet, different typefaces also provide different symbols. For example, and could theoretically appear in the same mathematical text with six different meanings. Normally, roman upright typeface is not used for symbols, except for symbols representing a standard function, such as the symbol "" of the sine function.
In order to have more symbols, and for allowing related mathematical objects to be represented by related symbols, diacritics, subscripts and superscripts are often used. For example, may denote the Fourier transform of the derivative of a function called
Other symbols
Symbols are not only used for naming mathematical objects. They can be used for operations for relations for logical connectives for quantifiers and for other purposes.
Some symbols are similar to Latin or Greek letters, some are obtained by deforming letters, some are traditional typographic symbols, but many have been specially designed for mathematics.
Expressions
An expression is a finite combination of symbols that is well-formed according to rules that depend on the context. In general, an expression denotes or names a mathematical object, and plays therefore in the language of mathematics the role of a noun phrase in the natural language.
An expression contains often some operators, and may therefore be evaluated by the action of the operators in it. For example, is an expression in which the operator can be evaluated for giving the result So, and are two different expressions that represent the same number. This is the meaning of the equality
A more complicated example is given by the expression that can be evaluated to Although the resulting expression contains the operators of division, subtraction and exponentiation, it cannot be evaluated further because and denote unspecified numbers.
History
Numbers
It is believed that a notation to represent numbers was first developed at least 50,000 years ago. Early mathematical ideas such as finger counting have also been represented by collections of rocks, sticks, bone, clay, stone, wood carvings, and knotted ropes. The tally stick is a way of counting dating back to the Upper Paleolithic. Perhaps the oldest known mathematical texts are those of ancient Sumer. The Census Quipu of the Andes and the Ishango Bone from Africa both used the tally mark method of accounting for numerical concepts.
The concept of zero and the introduction of a notation for it are important developments in early mathematics, which predates for centuries the concept of zero as a number. It was used as a placeholder by the Babylonians and Greek Egyptians, and then as an integer by the Mayans, Indians and Arabs (see the history of zero).
Modern notation
Until the 16th century, mathematics was essentially rhetorical, in the sense that everything but explicit numbers was expressed in words. However, some authors such as Diophantus used some symbols as abbreviations.
The first systematic use of formulas, and, in particular the use of symbols (variables) for unspecified numbers is generally attributed to François Viète (16th century). However, he used different symbols than those that are now standard.
Later, René Descartes (17th century) introduced the modern notation for variables and equations; in particular, the use of for unknown quantities and for known ones (constants). He introduced also the notation and the term "imaginary" for the imaginary unit.
The 18th and 19th centuries saw the standardization of mathematical notation as used today. Leonhard Euler was responsible for many of the notations currently in use: the functional notation for the base of the natural logarithm, for summation, etc. He also popularized the use of for the Archimedes constant (proposed by William Jones, based on an earlier notation of William Oughtred).
Since then many new notations have been introduced, often specific to a particular area of mathematics. Some notations are named after their inventors, such as Leibniz's notation, Legendre symbol, the Einstein summation convention, etc.
Typesetting
General typesetting systems are generally not well suited for mathematical notation. One of the reasons is that, in mathematical notation, the symbols are often arranged in two-dimensional figures, such as in:
TeX is a mathematically oriented typesetting system that was created in 1978 by Donald Knuth. It is widely used in mathematics, through its extension called LaTeX, and is a de facto standard. (The above expression is written in LaTeX.)
More recently, another approach for mathematical typesetting is provided by MathML. However, it is not well supported in web browsers, which is its primary target.
International standard mathematical notation
The international standard ISO 80000-2 (previously, ISO 31-11) specifies symbols for use in mathematical equations. The standard requires use of italic fonts for variables (e.g., ) and roman (upright) fonts for mathematical constants (e.g., e or π).
Non-Latin-based mathematical notation
Modern Arabic mathematical notation is based mostly on the Arabic alphabet and is used widely in the Arab world, especially in pre-tertiary education. (Western notation uses Arabic numerals, but the Arabic notation also replaces Latin letters and related symbols with Arabic script.)
In addition to Arabic notation, mathematics also makes use of Greek letters to denote a wide variety of mathematical objects and variables. On some occasions, certain Hebrew letters are also used (such as in the context of infinite cardinals).
Some mathematical notations are mostly diagrammatic, and so are almost entirely script independent. Examples are Penrose graphical notation and Coxeter–Dynkin diagrams.
Braille-based mathematical notations used by blind people include Nemeth Braille and GS8 Braille.
| Mathematics | Basics | null |
277197 | https://en.wikipedia.org/wiki/Estrildidae | Estrildidae | Estrildidae, or estrildid finches, is a family of small seed-eating passerine birds of the Old World tropics and Australasia. They comprise species commonly known as munias, mannikins, firefinches, parrotfinches and waxbills. Despite the word "finch" being included in the common names of some species, they are not closely related to birds with this name in other families, such as the Fringillidae, Emberizidae or Passerellidae.
They are gregarious and often colonial seed eaters with short, thick, but pointed bills. They are all similar in structure and habits, but vary widely in plumage colours and patterns.
All estrildids build large, domed nests and lay five to ten white eggs. Many species build roost nests. Some of the firefinches and pytilias are hosts to the brood-parasitic indigobirds and whydahs, respectively.
Most are sensitive to cold and require warm, usually tropical, habitats, although a few, such as the eastern alpine mannikin, mountain firetail, red-browed finch, and the genus Stagonopleura, have adapted to the cooler climates of southern Australia and the highlands of New Guinea.
The smallest species of the family is the Shelley's oliveback (Nesocharis shelleyi) at a mere , although the lightest species is the black-rumped waxbill (Estrilda troglodytes) at . The largest species is the Java sparrow (Padda oryzivora) at and .
Taxonomy
The family Estrildidae was introduced in 1850 by the French naturalist Charles Lucien Bonaparte as "Estreldinae", a spelling variant of the subfamily name. In the list of world birds maintained by Frank Gill, Pamela Rasmussen and David Donsker on behalf of the International Ornithological Committee (IOC) the family contains 140 species divided into 41 genera. Molecular phylogenetic studies have shown the family Estrildidae is sister to the family Viduidae containing the indigobirds and whydahs. The two families diverged around 15.5 million year ago. The most recent common ancestor of the Estrildidae is estimated to have lived around 10.9 million years ago. A genetic study of the Estrildidae by Urban Olsson and Per Alström published in 2020 identified 6 major clades. The radiations within these clades occurred between 4.5 and 8.9 million years ago. The authors proposed that each of these clades should be treated as a subfamily. This contrasts with an earlier proposal in which the family was divided into three subfamilies.
Genera list
| Biology and health sciences | Passerida | Animals |
277237 | https://en.wikipedia.org/wiki/Pumped-storage%20hydroelectricity | Pumped-storage hydroelectricity | Pumped-storage hydroelectricity (PSH), or pumped hydroelectric energy storage (PHES), is a type of hydroelectric energy storage used by electric power systems for load balancing.
A PSH system stores energy in the form of gravitational potential energy of water, pumped from a lower elevation reservoir to a higher elevation. Low-cost surplus off-peak electric power is typically used to run the pumps. During periods of high electrical demand, the stored water is released through turbines to produce electric power.
Pumped-storage hydroelectricity allows energy from intermittent sources (such as solar, wind, and other renewables) or excess electricity from continuous base-load sources (such as coal or nuclear) to be saved for periods of higher demand.
The reservoirs used with pumped storage can be quite small, when contrasted with the lakes of conventional hydroelectric plants of similar power capacity, and generating periods are often less than half a day.
The round-trip efficiency of PSH varies between 70% and 80%. Although the losses of the pumping process make the plant a net consumer of energy overall, the system increases revenue by selling more electricity during periods of peak demand, when electricity prices are highest. If the upper lake collects significant rainfall, or is fed by a river, then the plant may be a net energy producer in the manner of a traditional hydroelectric plant.
Pumped storage is by far the largest-capacity form of grid energy storage available, and, , accounts for around 95% of all active storage installations worldwide, with a total installed throughput capacity of over 181 GW and as of 2020 a total installed storage capacity of over 1.6 TWh.
Basic principle
A pumped-storage hydroelectricity generally consists of two water reservoirs at different heights, connected with each other.
At times of low electrical demand, excess generation capacity is used to pump water into the upper reservoir.
When there is higher demand, water is released back into the lower reservoir through a turbine, generating electricity. Pumped storage plants usually use reversible turbine/generator assemblies, which can act both as a pump and as a turbine generator (usually Francis turbine designs).
Variable speed operation further optimizes the round trip efficiency in pumped hydro storage plants.
In micro-PSH applications, a group of pumps and Pump As Turbine (PAT) could be implemented respectively for pumping and generating phases.
The same pump could be used in both modes by changing rotational direction and speed: the operation point in pumping usually differs from the operation point in PAT mode.
Types
In closed-loop systems, pure pumped-storage plants store water in an upper reservoir with no natural inflows, while pump-back plants utilize a combination of pumped storage and conventional hydroelectric plants with an upper reservoir that is replenished in part by natural inflows from a stream or river. Plants that do not use pumped storage are referred to as conventional hydroelectric plants; conventional hydroelectric plants that have significant storage capacity may be able to play a similar role in the electrical grid as pumped storage if appropriately equipped.
Economic efficiency
Taking into account conversion losses and evaporation losses from the exposed water surface, energy recovery of 70–80% or more can be achieved. This technique is currently the most cost-effective means of storing large amounts of electrical energy, but capital costs and the necessity of appropriate geography are critical decision factors in selecting pumped-storage plant sites.
The relatively low energy density of pumped storage systems requires either large flows and/or large differences in height between reservoirs. The only way to store a significant amount of energy is by having a large body of water located relatively near, but as high as possible above, a second body of water. In some places this occurs naturally, in others one or both bodies of water were man-made. Projects in which both reservoirs are artificial and in which no natural inflows are involved with either reservoir are referred to as "closed loop" systems.
These systems may be economical because they flatten out load variations on the power grid, permitting thermal power stations such as coal-fired plants and nuclear power plants that provide base-load electricity to continue operating at peak efficiency, while reducing the need for "peaking" power plants that use the same fuels as many base-load thermal plants, gas and oil, but have been designed for flexibility rather than maximal efficiency. Hence pumped storage systems are crucial when coordinating large groups of heterogeneous generators. Capital costs for pumped-storage plants are relatively high, although this is somewhat mitigated by their proven long service life of decades - and in some cases over a century, which is three to five times longer than utility-scale batteries. When electricity prices become negative, pumped hydro operators may earn twice - when "buying" the electricity to pump the water to the upper reservoir at negative spot prices and again when selling the electricity at a later time when prices are high.
Along with energy management, pumped storage systems help stabilize electrical network frequency and provide reserve generation. Thermal plants are much less able to respond to sudden changes in electrical demand that potentially cause frequency and voltage instability. Pumped storage plants, like other hydroelectric plants, can respond to load changes within seconds.
The most important use for pumped storage has traditionally been to balance baseload powerplants, but they may also be used to abate the fluctuating output of intermittent energy sources. Pumped storage provides a load at times of high electricity output and low electricity demand, enabling additional system peak capacity. In certain jurisdictions, electricity prices may be close to zero or occasionally negative on occasions that there is more electrical generation available than there is load available to absorb it. Although at present this is rarely due to wind or solar power alone, increased use of such generation will increase the likelihood of those occurrences.
It is particularly likely that pumped storage will become especially important as a balance for very large-scale photovoltaic and wind generation. Increased long-distance transmission capacity combined with significant amounts of energy storage will be a crucial part of regulating any large-scale deployment of intermittent renewable power sources. The high non-firm renewable electricity penetration in some regions supplies 40% of annual output, but 60% may be reached before additional storage is necessary.
Small-scale facilities
Smaller pumped storage plants cannot achieve the same economies of scale as larger ones, but some do exist, including a recent 13 MW project in Germany. Shell Energy has proposed a 5 MW project in Washington State. Some have proposed small pumped storage plants in buildings, although these are not yet economical. Also, it is difficult to fit large reservoirs into the urban landscape (and the fluctuating water level may make them unsuitable for recreational use). Nevertheless, some authors defend the technological simplicity and security of water supply as important externalities.
Location requirements
The main requirement for PSH is hilly country. The global greenfield pumped hydro atlas lists more than 800,000 potential sites around the world with combined storage of 86 million GWh (equivalent to the effective storage in about 2 trillion electric vehicle batteries), which is about 100 times more than needed to support 100% renewable electricity. Most are closed-loop systems away from rivers. Areas of natural beauty and new dams on rivers can be avoided because of the very large number of potential sites. Some projects utilise existing reservoirs (dubbed "bluefield") such as the 350 Gigawatt-hour Snowy 2.0 scheme under construction in Australia. Some recently proposed projects propose to take advantage of "brownfield" locations such as disused mines such as the Kidston project under construction in Australia.
Environmental impact
Water requirements for PSH are small: about 1 gigalitre of initial fill water per gigawatt-hour of storage. This water is recycled uphill and back downhill between the two reservoirs for many decades, but evaporation losses (beyond what rainfall and any inflow from local waterways provide) must be replaced. Land requirements are also small: about 10 hectares per gigawatt-hour of storage, which is much smaller than the land occupied by the solar and windfarms that the storage might support. Closed loop (off-river) pumped hydro storage has the smallest carbon emissions per unit of storage of all candidates for large-scale energy storage.
Potential technologies
Seawater
Pumped storage plants can operate with seawater, although there are additional challenges compared to using fresh water, such as saltwater corrosion and barnacle growth. Inaugurated in 1966, the 240 MW Rance tidal power station in France can partially work as a pumped-storage station. When high tides occur at off-peak hours, the turbines can be used to pump more seawater into the reservoir than the high tide would have naturally brought in. It is the only large-scale power plant of its kind.
In 1999, the 30 MW Yanbaru project in Okinawa was the first demonstration of seawater pumped storage. It has since been decommissioned. A 300 MW seawater-based Lanai Pumped Storage Project was considered for Lanai, Hawaii, and seawater-based projects have been proposed in Ireland. A pair of proposed projects in the Atacama Desert in northern Chile would use 600 MW of photovoltaic solar (Skies of Tarapacá) together with 300 MW of pumped storage (Mirror of Tarapacá) lifting seawater up a coastal cliff.
Freshwater coastal reservoirs
Freshwater from the river floods is stored in the sea area replacing seawater by constructing coastal reservoirs. The stored river water is pumped to uplands by constructing a series of embankment canals and pumped storage hydroelectric stations for the purpose of energy storage, irrigation, industrial, municipal, rejuvenation of over exploited rivers, etc. These multipurpose coastal reservoir projects offer massive pumped-storage hydroelectric potential to utilize variable and intermittent solar and wind power that are carbon-neutral, clean, and renewable energy sources.
Underground reservoirs
The use of underground reservoirs has been investigated. Recent examples include the proposed Summit project in Norton, Ohio, the proposed Maysville project in Kentucky (underground limestone mine), and the Mount Hope project in New Jersey, which was to have used a former iron mine as the lower reservoir. The proposed energy storage at the Callio site in Pyhäjärvi (Finland) would utilize the deepest base metal mine in Europe, with elevation difference. Several new underground pumped storage projects have been proposed. Cost-per-kilowatt estimates for these projects can be lower than for surface projects if they use existing underground mine space. There are limited opportunities involving suitable underground space, but the number of underground pumped storage opportunities may increase if abandoned coal mines prove suitable.
In Bendigo, Victoria, Australia, the Bendigo Sustainability Group has proposed the use of the old gold mines under Bendigo for Pumped Hydro Energy Storage. Bendigo has the greatest concentration of deep shaft hard rock mines anywhere in the world with over 5,000 shafts sunk under Bendigo in the second half of the 19th Century. The deepest shaft extends 1,406 metres vertically underground. A recent pre-feasibility study has shown the concept to be viable with a generation capacity of 30 MW and a run time of 6 hours using a water head of over 750 metres.
US-based start-up Quidnet Energy is exploring using abandoned oil and gas wells for pumped storage. If successful they hope to scale up, utilizing some of the 3 million abandoned wells in the US.
Using hydraulic fracturing pressure can be stored underground in impermeable strata such as shale. The shale used contains no hydrocarbons.
Decentralised systems
Small (or micro) applications for pumped storage could be built on streams and within infrastructures, such as drinking water networks and artificial snow-making infrastructures. In this regard, a storm-water basin has been concretely implemented as a cost-effective solution for a water reservoir in a micro-pumped hydro energy storage. Such plants provide distributed energy storage and distributed flexible electricity production and can contribute to the decentralized integration of intermittent renewable energy technologies, such as wind power and solar power. Reservoirs that can be used for small pumped-storage hydropower plants could include natural or artificial lakes, reservoirs within other structures such as irrigation, or unused portions of mines or underground military installations. In Switzerland one study suggested that the total installed capacity of small pumped-storage hydropower plants in 2011 could be increased by 3 to 9 times by providing adequate policy instruments.
Using a pumped-storage system of cisterns and small generators, pico hydro may also be effective for "closed loop" home energy generation systems.
Underwater reservoirs
In March 2017, the research project StEnSea (Storing Energy at Sea) announced their successful completion of a four-week test of a pumped storage underwater reservoir. In this configuration, a hollow sphere submerged and anchored at great depth acts as the lower reservoir, while the upper reservoir is the enclosing body of water. Electricity is created when water is let in via a reversible turbine integrated into the sphere. During off-peak hours, the turbine changes direction and pumps the water out again, using "surplus" electricity from the grid.
The quantity of power created when water is let in, grows proportionally to the height of the column of water above the sphere. In other words: the deeper the sphere is located, the more densely it can store energy.
As such, the energy storage capacity of the submerged reservoir is not governed by the gravitational energy in the traditional sense, but by the vertical pressure variation.
High-density pumped hydro
RheEnergise aim to improve the efficiency of pumped storage by using fluid 2.5x denser than water ("a fine-milled suspended solid in water"), such that "projects can be 2.5x smaller for the same power."
History
The first use of pumped storage was in 1907 in Switzerland, at the Engeweiher pumped storage facility near Schaffhausen, Switzerland. In the 1930s reversible hydroelectric turbines became available. This apparatus could operate both as turbine generators and in reverse as electric motor-driven pumps. The latest in large-scale engineering technology is variable speed machines for greater efficiency. These machines operate in synchronization with the network frequency when generating, but operate asynchronously (independent of the network frequency) when pumping.
The first use of pumped-storage in the United States was in 1930 by the Connecticut Electric and Power Company, using a large reservoir located near New Milford, Connecticut, pumping water from the Housatonic River to the storage reservoir above.
Worldwide use
In 2009, world pumped storage generating capacity was 104 GW, while other sources claim 127 GW, which comprises the vast majority of all types of utility grade electric storage.
The European Union had 38.3 GW net capacity (36.8% of world capacity) out of a total of 140 GW of hydropower and representing 5% of total net electrical capacity in the EU.
Japan had 25.5 GW net capacity (24.5% of world capacity).
The six largest operational pumped-storage plants are listed below (for a detailed list see List of pumped-storage hydroelectric power stations):
Australia
Australia has 15GW of pumped storage under construction or in development. Examples include:
In June 2018 the Australian federal government announced that 14 sites had been identified in Tasmania for pumped storage hydro, with the potential of adding 4.8GW to the national grid if a second interconnector beneath Bass Strait was constructed.
The Snowy 2.0 project will link two existing dams in the New South Wales' Snowy Mountains to provide 2,000 MW of capacity and 350,000 MWh of storage.
In September 2022, a pumped hydroelectric storage (PHES) scheme was announced at Pioneer-Burdekin in central Queensland that has the potential to be the largest PHES in the world at 5 GW.
China
China has the largest capacity of pumped-storage hydroelectricity in the world.
In January 2019, the State Grid Corporation of China announced plans to invest US$5.7 billion in five pumped hydro storage plants with a total 6 GW capacity, to be located in Hebei, Jilin, Zhejiang, Shandong provinces, and in Xinjiang Autonomous Region. China is seeking to build 40 GW of pumped hydro capacity installed by 2020.
Norway
There are 9 power stations capable of pumping with a total installed capacity of 1344 MW and an average annual production of 2247 GWh. The pumped storage hydropower in Norway is built a bit differently from the rest of the world. They are designed for seasonal pumping. Most of them can also not cycle the water endlessly, but only pump and reuse once. The reason for this is the design of the tunnels and the elevation of lower and upper reservoirs. Some, like Nygard power station, pump water from several river intakes up to a reservoir.
The largest one, Saurdal, which is part of the Ulla-Førre complex, has four 160 MW Francis turbines, but only two are reversible. The lower reservoir is at a higher elevation than the station itself, and thus the water pumped up can only be used once before it has to flow to the next station, Kvilldal, further down the tunnel system. And in addition to the lower reservoir, it will receive water that can be pumped up from 23 river/stream and small reservoir intakes. Some of which will have already gone through a smaller power station on its way.
United States
In 2010, the United States had 21.5 GW of pumped storage generating capacity (20.6% of world capacity).
PSH contributed 21,073 GWh of energy in 2020 in the United States, but −5,321 GWh (net) because more energy is consumed in pumping than is generated. Nameplate pumped storage capacity had grown to 21.6 GW by 2014, with pumped storage comprising 97% of grid-scale energy storage in the United States. As of late 2014, there were 51 active project proposals with a total of 39 GW of new nameplate capacity across all stages of the FERC licensing process for new pumped storage hydroelectric plants in the United States, but no new plants were currently under construction in the United States at the time.
Italy
Italy reached peak usage of pumped storage (pompaggi) in 2003, with about 8 TWh. For decades, Italy had excess capacity because its own nuclear program was interrupted in the 1980s, so pumping stations are mostly operated by night when France exports surplus nuclear electricity at near-zero prices. In 2019, the grid operator wanted 6 GW of extra capacity to be built in central and Southern Italy. In 2024, Edison planned 500 MW new capacity.
Hybrid systems
Conventional hydroelectric dams may also make use of pumped storage in a hybrid system that both generates power from water naturally flowing into the reservoir as well as storing water pumped back to the reservoir from below the dam. The Grand Coulee Dam in the United States was expanded with a pump-back system in 1973. Existing dams may be repowered with reversing turbines thereby extending the length of time the plant can operate at capacity. Optionally a pump back powerhouse such as the Russell Dam (1992) may be added to a dam for increased generating capacity. Making use of an existing dam's upper reservoir and transmission system can expedite projects and reduce costs.
| Technology | Energy storage | null |
277248 | https://en.wikipedia.org/wiki/Blinded%20experiment | Blinded experiment | In a blind or blinded experiment, information which may influence the participants of the experiment is withheld until after the experiment is complete. Good blinding can reduce or eliminate experimental biases that arise from a participants' expectations, observer's effect on the participants, observer bias, confirmation bias, and other sources. A blind can be imposed on any participant of an experiment, including subjects, researchers, technicians, data analysts, and evaluators. In some cases, while blinding would be useful, it is impossible or unethical. For example, it is not possible to blind a patient to their treatment in a physical therapy intervention. A good clinical protocol ensures that blinding is as effective as possible within ethical and practical constraints.
During the course of an experiment, a participant becomes unblinded if they deduce or otherwise obtain information that has been masked to them. For example, a patient who experiences a side effect may correctly guess their treatment, becoming unblinded. Unblinding is common in blinded experiments, particularly in pharmacological trials. In particular, trials on pain medication and antidepressants are poorly blinded. Unblinding that occurs before the conclusion of a study is a source of experimental error, as the bias that was eliminated by blinding is re-introduced. The CONSORT reporting guidelines recommend that all studies assess and report unblinding. In practice, very few studies do so.
Blinding is an important tool of the scientific method, and is used in many fields of research. In some fields, such as medicine, it is considered essential. In clinical research, a trial that is not a blinded trial is called an open trial.
History
The first known blind experiment was conducted by the French Royal Commission on Animal Magnetism in 1784 to investigate the claims of mesmerism as proposed by Charles d'Eslon, a former associate of Franz Mesmer. In the investigations, the researchers (physically) blindfolded mesmerists and asked them to identify objects that the experimenters had previously filled with "vital fluid". The subjects were unable to do so.
In 1817, the first blind experiment recorded to have occurred outside of a scientific setting compared the musical quality of a Stradivarius violin to one with a guitar-like design. A violinist played each instrument while a committee of scientists and musicians listened from another room so as to avoid prejudice.
An early example of a double-blind protocol was the Nuremberg salt test of 1835 performed by Friedrich Wilhelm von Hoven, Nuremberg's highest-ranking public health official, as well as a close friend of Friedrich Schiller. This trial contested the effectiveness of homeopathic dilution.
In 1865, Claude Bernard published his Introduction to the Study of Experimental Medicine, which advocated for the blinding of researchers. Bernard's recommendation that an experiment's observer should not know the hypothesis being tested contrasted starkly with the prevalent Enlightenment-era attitude that scientific observation can only be objectively valid when undertaken by a well-educated, informed scientist. The first study recorded to have a blinded researcher was conducted in 1907 by W. H. R. Rivers and H. N. Webber to investigate the effects of caffeine. The need to blind researchers became widely recognized in the mid-20th century.
Background
Bias
A number of biases are present when a study is insufficiently blinded. Patient-reported outcomes can be different if the patient is not blinded to their treatment. Likewise, failure to blind researchers results in observer bias. Unblinded data analysts may favor an analysis that supports their existing beliefs (confirmation bias). These biases are typically the result of subconscious influences, and are present even when study participants believe they are not influenced by them.
Terminology
In medical research, the terms single-blind, double-blind and triple-blind are commonly used to describe blinding. These terms describe experiments in which (respectively) one, two, or three parties are blinded to some information. Most often, single-blind studies blind patients to their treatment allocation, double-blind studies blind both patients and researchers to treatment allocations, and triple-blinded studies blind patients, researcher, and some other third party (such as a monitoring committee) to treatment allocations. However, the meaning of these terms can vary from study to study.
CONSORT guidelines state that these terms should no longer be used because they are ambiguous. For instance, "double-blind" could mean that the data analysts and patients were blinded; or the patients and outcome assessors were blinded; or the patients and people offering the intervention were blinded, etc. The terms also fail to convey the information that was masked and the amount of unblinding that occurred. It is not sufficient to specify the number of parties that have been blinded. To describe an experiment's blinding, it is necessary to report who has been blinded to what information, and how well each blind succeeded.
Unblinding
"Unblinding" occurs in a blinded experiment when information becomes available to one from whom it has been masked. In clinical studies, unblinding may occur unintentionally when a patient deduces their treatment group. Unblinding that occurs before the conclusion of an experiment is a source of bias. Some degree of premature unblinding is common in blinded experiments. When a blind is imperfect, its success is judged on a spectrum with no blind (or complete failure of blinding) on one end, perfect blinding on the other, and poor or good blinding between. Thus, the common view of studies as blinded or unblinded is an example of a false dichotomy.
Success of blinding is assessed by questioning study participants about information that has been masked to them (e.g. did the participant receive the drug or placebo?). In a perfectly blinded experiment, the responses should be consistent with no knowledge of the masked information. However, if unblinding has occurred, the responses will indicate the degree of unblinding. Since unblinding cannot be measured directly, but must be inferred from participants' responses, its measured value will depend on the nature of the questions asked. As a result, it is not possible to measure unblinding in a way that is completely objective. Nonetheless, it is still possible to make informed judgments about the quality of a blind. Poorly blinded studies rank above unblinded studies and below well-blinded studies in the hierarchy of evidence.
Post-study unblinding
Post-study unblinding is the release of masked data upon completion of a study. In clinical studies, post-study unblinding serves to inform subjects of their treatment allocation. Removing a blind upon completion of a study is never mandatory, but is typically performed as a courtesy to study participants. Unblinding that occurs after the conclusion of a study is not a source of bias, because data collection and analysis are both complete at this time.
Premature unblinding
Premature unblinding is any unblinding that occurs before the conclusion of a study. In contrast with post-study unblinding, premature unblinding is a source of bias. A code-break procedure dictates when a subject should be unblinded prematurely. A code-break procedure should only allow for unblinding in cases of emergency. Unblinding that occurs in compliance with code-break procedure is strictly documented and reported.
Premature unblinding may also occur when a participant infers from experimental conditions information that has been masked to them. A common cause for unblinding is the presence of side effects (or effects) in the treatment group. In pharmacological trials, premature unblinding can be reduced with the use of an active placebo, which conceals treatment allocation by ensuring the presence of side effects in both groups. However, side effects are not the only cause of unblinding; any perceptible difference between the treatment and control groups can contribute to premature unblinding.
A problem arises in the assessment of blinding because asking subjects to guess masked information may prompt them to try to infer that information. Researchers speculate that this may contribute to premature unblinding. Furthermore, it has been reported that some subjects of clinical trials attempt to determine if they have received an active treatment by gathering information on social media and message boards. While researchers counsel patients not to use social media to discuss clinical trials, their accounts are not monitored. This behavior is believed to be a source of unblinding. CONSORT standards and good clinical practice guidelines recommend the reporting of all premature unblinding. In practice, unintentional unblinding is rarely reported.
Significance
Bias due to poor blinding tends to favor the experimental group, resulting in inflated effect size and risk of false positives. Success or failure of blinding is rarely reported or measured; it is implicitly assumed that experiments reported as "blind" are truly blind. Critics have pointed out that without assessment and reporting, there is no way to know if a blind succeeded. This shortcoming is especially concerning given that even a small error in blinding can produce a statistically significant result in the absence of any real difference between test groups when a study is sufficiently powered (i.e. statistical significance is not robust to bias). As such, many statistically significant results in randomized controlled trials may be caused by error in blinding. Some researchers have called for the mandatory assessment of blinding efficacy in clinical trials.
Applications
In medicine
Blinding is considered essential in medicine, but is often difficult to achieve. For example, it is difficult to compare surgical and non-surgical interventions in blind trials. In some cases, sham surgery may be necessary for the blinding process. A good clinical protocol ensures that blinding is as effective as possible within ethical and practical constrains.
Studies of blinded pharmacological trials across widely varying domains find evidence of high levels of unblinding. Unblinding has been shown to affect both patients and clinicians. This evidence challenges the common assumption that blinding is highly effective in pharmacological trials. Unblinding has also been documented in clinical trials outside of pharmacology.
Pain
A 2018 meta-analysis found that assessment of blinding was reported in only 23 out of 408 randomized controlled trials for chronic pain (5.6%). The study concluded upon analysis of pooled data that the overall quality of the blinding was poor, and the blinding was "not successful." Additionally, both pharmaceutical sponsorship and the presence of side effects were associated with lower rates of reporting assessment of blinding.
Depression
Studies have found evidence of extensive unblinding in antidepressant trials: at least three-quarters of patients were able to correctly guess their treatment assignment. Unblinding also occurs in clinicians. Better blinding of patients and clinicians reduces effect size. Researchers concluded that unblinding inflates effect size in antidepressant trials. Some researchers believe that antidepressants are not effective for the treatment of depression and only outperform placebos due to systematic error. These researchers argue that antidepressants are just active placebos.
Acupuncture
While the possibility of blinded trials on acupuncture is controversial, a 2003 review of 47 randomized controlled trials found no fewer than four methods of blinding patients to acupuncture treatment: 1) superficial needling of true acupuncture points, 2) use of acupuncture points which are not indicated for the condition being treated, 3) insertion of needles outside of true acupuncture points, and 4) the use of placebo needles which are designed not to penetrate the skin. The authors concluded that there was "no clear association between type of sham intervention used and the results of the trials."
A 2018 study on acupuncture which used needles that did not penetrate the skin as a sham treatment found that 68% of patients and 83% of acupuncturists correctly identified their group allocation. The authors concluded that the blinding had failed, but that more advanced placebos may someday offer the possibility of well-blinded studies in acupuncture.
In physics
It is standard practice in physics to perform blinded data analysis. After data analysis is complete, one is allowed to unblind the data. A prior agreement to publish the data regardless of the results of the analysis may be made to prevent publication bias.
In social sciences
Social science research is particularly prone to observer bias, so it is important in these fields to properly blind the researchers. In some cases, while blind experiments would be useful, they are impractical or unethical. Blinded data analysis can reduce bias, but is rarely used in social science research.
In forensics
In a police photo lineup, an officer shows a group of photos to a witness and asks the witness to identify the individual who committed the crime. Since the officer is typically aware of who the suspect is, they may (subconsciously or consciously) influence the witness to choose the individual that they believe committed the crime. There is a growing movement in law enforcement to move to a blind procedure in which the officer who shows the photos to the witness does not know who the suspect is.
In music
Auditions for symphony orchestras take place behind a curtain so that the judges cannot see the performer. Blinding the judges to the gender of the performers has been shown to increase the hiring of women. Blind tests can also be used to compare the quality of musical instruments.
| Physical sciences | Science basics | Basics and measurement |
277266 | https://en.wikipedia.org/wiki/Iron%20oxide | Iron oxide | An iron oxide is a chemical compound composed of iron and oxygen. Several iron oxides are recognized. Often they are non-stoichiometric. Ferric oxyhydroxides are a related class of compounds, perhaps the best known of which is rust.
Iron oxides and oxyhydroxides are widespread in nature and play an important role in many geological and biological processes. They are used as iron ores, pigments, catalysts, and in thermite, and occur in hemoglobin. Iron oxides are inexpensive and durable pigments in paints, coatings and colored concretes. Colors commonly available are in the "earthy" end of the yellow/orange/red/brown/black range. When used as a food coloring, it has E number E172.
Stoichiometries
Iron oxides feature as ferrous (Fe(II)) or ferric (Fe(III)) or both. They adopt octahedral or tetrahedral coordination geometry. Only a few oxides are significant at the earth's surface, particularly wüstite, magnetite, and hematite.
Oxides of FeII
FeO: iron(II) oxide, wüstite
Mixed oxides of FeII and FeIII
Fe3O4: Iron(II,III) oxide, magnetite
Fe4O5
Fe5O6
Fe5O7
Fe25O32
Fe13O19
Oxides of FeIII
Fe2O3: iron(III) oxide
α-Fe2O3: alpha phase, hematite
β-Fe2O3: beta phase
γ-Fe2O3: gamma phase, maghemite
ε-Fe2O3: epsilon phase
Thermal expansion
Oxide-hydroxides
goethite (α-FeOOH)
akaganéite (β-FeOOH)
lepidocrocite (γ-FeOOH)
feroxyhyte (δ-FeOOH)
ferrihydrite (Fe5HO8 · 4 H2O approx., or 5 Fe2O3 · 9 H2O, better recast as FeOOH · 0.4 H2O)
high-pressure pyrite-structured FeOOH. Once dehydration is triggered, this phase may form FeO2Hx (0 < x < 1).
green rust (FeFeOH3x + y − z (A−)z where A− is Cl− or 0.5 )
Reactions
In blast furnaces and related factories, iron oxides are converted to the metal. Typical reducing agents are various forms of carbon. A representative reaction starts with ferric oxide:
In nature
Iron is stored in many organisms in the form of ferritin, which is a ferrous oxide encased in a solubilizing protein sheath.
Species of bacteria, including Shewanella oneidensis, Geobacter sulfurreducens and Geobacter metallireducens, use iron oxides as terminal electron acceptors.
Uses
Almost all iron ores are oxides, so in that sense these materials are important precursors to iron metal and its many alloys.
Iron oxides are important pigments, coming in a variety of colors (black, red, yellow). Among their many advantages, they are inexpensive, strongly colored, and nontoxic.
Magnetite is a component of magnetic recording tapes.
| Physical sciences | Oxide salts | Chemistry |
277289 | https://en.wikipedia.org/wiki/Wind%20power | Wind power | Wind power is the use of wind energy to generate useful work. Historically, wind power was used by sails, windmills and windpumps, but today it is mostly used to generate electricity. This article deals only with wind power for electricity generation.
Today, wind power is generated almost completely with wind turbines, generally grouped into wind farms and connected to the electrical grid.
In 2022, wind supplied over 2,304 TWh of electricity, which was 7.8% of world electricity.
With about 100 GW added during 2021, mostly in China and the United States, global installed wind power capacity exceeded 800 GW. 32 countries generated more than a tenth of their electricity from wind power in 2023 and wind generation has nearly tripled since 2015. To help meet the Paris Agreement goals to limit climate change, analysts say it should expand much faster – by over 1% of electricity generation per year.
Wind power is considered a sustainable, renewable energy source, and has a much smaller impact on the environment compared to burning fossil fuels. Wind power is variable, so it needs energy storage or other dispatchable generation energy sources to attain a reliable supply of electricity. Land-based (onshore) wind farms have a greater visual impact on the landscape than most other power stations per energy produced. Wind farms sited offshore have less visual impact and have higher capacity factors, although they are generally more expensive. Offshore wind power currently has a share of about 10% of new installations.
Wind power is one of the lowest-cost electricity sources per unit of energy produced.
In many locations, new onshore wind farms are cheaper than new coal or gas plants.
Regions in the higher northern and southern latitudes have the highest potential for wind power. In most regions, wind power generation is higher in nighttime, and in winter when solar power output is low. For this reason, combinations of wind and solar power are suitable in many countries.
Wind energy resources
Wind is air movement in the Earth's atmosphere. In a unit of time, say 1 second, the volume of air that had passed an area is . If the air density is , the mass of this volume of air is , and the power transfer, or energy transfer per second is . Wind power is thus proportional to the third power of the wind speed; the available power increases eightfold when the wind speed doubles. Change of wind speed by a factor of 2.1544 increases the wind power by one order of magnitude (multiply by 10).
The global wind kinetic energy averaged approximately 1.50 MJ/m2 over the period from 1979 to 2010, 1.31 MJ/m2 in the Northern Hemisphere with 1.70 MJ/m2 in the Southern Hemisphere. The atmosphere acts as a thermal engine, absorbing heat at higher temperatures, releasing heat at lower temperatures. The process is responsible for the production of wind kinetic energy at a rate of 2.46 W/m2 thus sustaining the circulation of the atmosphere against friction.
Through wind resource assessment, it is possible to estimate wind power potential globally, by country or region, or for a specific site. The Global Wind Atlas provided by the Technical University of Denmark in partnership with the World Bank provides a global assessment of wind power potential.
Unlike 'static' wind resource atlases which average estimates of wind speed and power density across multiple years, tools such as Renewables.ninja provide time-varying simulations of wind speed and power output from different wind turbine models at an hourly resolution. More detailed, site-specific assessments of wind resource potential can be obtained from specialist commercial providers, and many of the larger wind developers have in-house modeling capabilities.
The total amount of economically extractable power available from the wind is considerably more than present human power use from all sources. The strength of wind varies, and an average value for a given location does not alone indicate the amount of energy a wind turbine could produce there.
To assess prospective wind power sites, a probability distribution function is often fit to the observed wind speed data. Different locations will have different wind speed distributions. The Weibull model closely mirrors the actual distribution of hourly/ten-minute wind speeds at many locations. The Weibull factor is often close to 2 and therefore a Rayleigh distribution can be used as a less accurate, but simpler model.
Wind farms
A wind farm is a group of wind turbines in the same location. A large wind farm may consist of several hundred individual wind turbines distributed over an extended area. The land between the turbines may be used for agricultural or other purposes. A wind farm may also be located offshore. Almost all large wind turbines have the same design — a horizontal axis wind turbine having an upwind rotor with 3 blades, attached to a nacelle on top of a tall tubular tower.
In a wind farm, individual turbines are interconnected with a medium voltage (often 34.5 kV) power collection system and communications network. In general, a distance of 7D (7 times the rotor diameter of the wind turbine) is set between each turbine in a fully developed wind farm. At a substation, this medium-voltage electric current is increased in voltage with a transformer for connection to the high voltage electric power transmission system.
Generator characteristics and stability
Most modern turbines use variable speed generators combined with either a partial or full-scale power converter between the turbine generator and the collector system, which generally have more desirable properties for grid interconnection and have low voltage ride through-capabilities. Modern turbines use either doubly fed electric machines with partial-scale converters or squirrel-cage induction generators or synchronous generators (both permanently and electrically excited) with full-scale converters. Black start is possible and is being further developed for places (such as Iowa) which generate most of their electricity from wind.
Transmission system operators will supply a wind farm developer with a grid code to specify the requirements for interconnection to the transmission grid. This will include the power factor, the constancy of frequency, and the dynamic behaviour of the wind farm turbines during a system fault.
Offshore wind power
Offshore wind power is wind farms in large bodies of water, usually the sea. These installations can use the more frequent and powerful winds that are available in these locations and have less visual impact on the landscape than land-based projects. However, the construction and maintenance costs are considerably higher.
As of November 2021, the Hornsea Wind Farm in the United Kingdom is the largest offshore wind farm in the world at 1,218 MW.
Collection and transmission network
Near offshore wind farms may be connected by AC and far offshore by HVDC.
Wind power resources are not always located near to high population density. As transmission lines become longer, the losses associated with power transmission increase, as modes of losses at lower lengths are exacerbated and new modes of losses are no longer negligible as the length is increased; making it harder to transport large loads over large distances.
When the transmission capacity does not meet the generation capacity, wind farms are forced to produce below their full potential or stop running altogether, in a process known as curtailment. While this leads to potential renewable generation left untapped, it prevents possible grid overload or risk to reliable service.
One of the biggest current challenges to wind power grid integration in some countries is the necessity of developing new transmission lines to carry power from wind farms, usually in remote lowly populated areas due to availability of wind, to high load locations, usually on the coasts where population density is higher. Any existing transmission lines in remote locations may not have been designed for the transport of large amounts of energy. In particular geographic regions, peak wind speeds may not coincide with peak demand for electrical power, whether offshore or onshore. A possible future option may be to interconnect widely dispersed geographic areas with an HVDC super grid.
Wind power capacity and production
Growth trends
In 2020, wind supplied almost 1600 TWh of electricity, which was over 5% of worldwide electrical generation and about 2% of energy consumption. With over 100 GW added during 2020, mostly in China, global installed wind power capacity reached more than 730 GW. But to help meet the Paris Agreement's goals to limit climate change, analysts say it should expand much faster – by over 1% of electricity generation per year. Expansion of wind power is being hindered by fossil fuel subsidies.
The actual amount of electric power that wind can generate is calculated by multiplying the nameplate capacity by the capacity factor, which varies according to equipment and location. Estimates of the capacity factors for wind installations are in the range of 35% to 44%.
Capacity factor
Since wind speed is not constant, a wind farm's annual energy production is never as much as the sum of the generator nameplate ratings multiplied by the total hours in a year. The ratio of actual productivity in a year to this theoretical maximum is called the capacity factor. Online data is available for some locations, and the capacity factor can be calculated from the yearly output.
Penetration
Wind energy penetration is the fraction of energy produced by wind compared with the total generation. Wind power's share of worldwide electricity usage in 2021 was almost 7%, up from 3.5% in 2015.
There is no generally accepted maximum level of wind penetration. The limit for a particular grid will depend on the existing generating plants, pricing mechanisms, capacity for energy storage, demand management, and other factors. An interconnected electric power grid will already include reserve generating and transmission capacity to allow for equipment failures. This reserve capacity can also serve to compensate for the varying power generation produced by wind stations. Studies have indicated that 20% of the total annual electrical energy consumption may be incorporated with minimal difficulty. These studies have been for locations with geographically dispersed wind farms, some degree of dispatchable energy or hydropower with storage capacity, demand management, and interconnected to a large grid area enabling the export of electric power when needed. Electrical utilities continue to study the effects of large-scale penetration of wind generation on system stability.
A wind energy penetration figure can be specified for different duration of time but is often quoted annually. To generate almost all electricity from wind annually requires substantial interconnection to other systems, for example some wind power in Scotland is sent to the rest of the British grid. On a monthly, weekly, daily, or hourly basis—or less—wind might supply as much as or more than 100% of current use, with the rest stored, exported or curtailed. The seasonal industry might then take advantage of high wind and low usage times such as at night when wind output can exceed normal demand. Such industry might include the production of silicon, aluminum, steel, or natural gas, and hydrogen, and using future long-term storage to facilitate 100% energy from variable renewable energy. Homes and businesses can also be programmed to vary electricity demand, for example by remotely turning up water heater thermostats.
Variability
Wind power is variable, and during low wind periods, it may need to be replaced by other power sources. Transmission networks presently cope with outages of other generation plants and daily changes in electrical demand, but the variability of intermittent power sources such as wind power is more frequent than those of conventional power generation plants which, when scheduled to be operating, may be able to deliver their nameplate capacity around 95% of the time.
Electric power generated from wind power can be highly variable at several different timescales: hourly, daily, or seasonally. Annual variation also exists but is not as significant. Because instantaneous electrical generation and consumption must remain in balance to maintain grid stability, this variability can present substantial challenges to incorporating large amounts of wind power into a grid system. Intermittency and the non-dispatchable nature of wind energy production can raise costs for regulation, incremental operating reserve, and (at high penetration levels) could require an increase in the already existing energy demand management, load shedding, storage solutions, or system interconnection with HVDC cables.
Fluctuations in load and allowance for the failure of large fossil-fuel generating units require operating reserve capacity, which can be increased to compensate for the variability of wind generation.
Utility-scale batteries are often used to balance hourly and shorter timescale variation, but car batteries may gain ground from the mid-2020s. Wind power advocates argue that periods of low wind can be dealt with by simply restarting existing power stations that have been held in readiness, or interlinking with HVDC.
The combination of diversifying variable renewables by type and location, forecasting their variation, and integrating them with dispatchable renewables, flexible fueled generators, and demand response can create a power system that has the potential to meet power supply needs reliably. Integrating ever-higher levels of renewables is being successfully demonstrated in the real world.
Solar power tends to be complementary to wind. On daily to weekly timescales, high-pressure areas tend to bring clear skies and low surface winds, whereas low-pressure areas tend to be windier and cloudier. On seasonal timescales, solar energy peaks in summer, whereas in many areas wind energy is lower in summer and higher in winter. Thus the seasonal variation of wind and solar power tend to cancel each other somewhat. Wind hybrid power systems are becoming more popular.
Predictability
For any particular generator, there is an 80% chance that wind output will change less than 10% in an hour and a 40% chance that it will change 10% or more in 5 hours.
In summer 2021, wind power in the United Kingdom fell due to the lowest winds in seventy years, In the future, smoothing peaks by producing green hydrogen may help when wind has a larger share of generation.
While the output from a single turbine can vary greatly and rapidly as local wind speeds vary, as more turbines are connected over larger and larger areas the average power output becomes less variable and more predictable. Weather forecasting permits the electric-power network to be readied for the predictable variations in production that occur.
It is thought that the most reliable low-carbon electricity systems will include a large share of wind power.
Energy storage
Typically, conventional hydroelectricity complements wind power very well. When the wind is blowing strongly, nearby hydroelectric stations can temporarily hold back their water. When the wind drops they can, provided they have the generation capacity, rapidly increase production to compensate. This gives a very even overall power supply and virtually no loss of energy and uses no more water.
Alternatively, where a suitable head of water is not available, pumped-storage hydroelectricity or other forms of grid energy storage such as compressed air energy storage and thermal energy storage can store energy developed by high-wind periods and release it when needed. The type of storage needed depends on the wind penetration level – low penetration requires daily storage, and high penetration requires both short- and long-term storage – as long as a month or more. Stored energy increases the economic value of wind energy since it can be shifted to displace higher-cost generation during peak demand periods. The potential revenue from this arbitrage can offset the cost and losses of storage. Although pumped-storage power systems are only about 75% efficient and have high installation costs, their low running costs and ability to reduce the required electrical base-load can save both fuel and total electrical generation costs.
Energy payback
The energy needed to build a wind farm divided into the total output over its life, Energy Return on Energy Invested, of wind power varies, but averages about 20–25. Thus, the energy payback time is typically around a year.
Economics
Onshore wind is an inexpensive source of electric power, cheaper than coal plants and new gas plants. According to BusinessGreen, wind turbines reached grid parity (the point at which the cost of wind power matches traditional sources) in some areas of Europe in the mid-2000s, and in the US around the same time. Falling prices continue to drive the Levelized cost down and it has been suggested that it has reached general grid parity in Europe in 2010, and will reach the same point in the US around 2016 due to an expected reduction in capital costs of about 12%. In 2021, the CEO of Siemens Gamesa warned that increased demand for low-cost wind turbines combined with high input costs and high costs of steel result in increased pressure on the manufacturers and decreasing profit margins.
Northern Eurasia, Canada, some parts of the United States, and Patagonia in Argentina are the best areas for onshore wind: whereas in other parts of the world solar power, or a combination of wind and solar, tend to be cheaper.
Electric power cost and trends
Wind power is capital intensive but has no fuel costs. The price of wind power is therefore much more stable than the volatile prices of fossil fuel sources. However, the estimated average cost per unit of electric power must incorporate the cost of construction of the turbine and transmission facilities, borrowed funds, return to investors (including the cost of risk), estimated annual production, and other components, averaged over the projected useful life of the equipment, which may be more than 20 years. Energy cost estimates are highly dependent on these assumptions so published cost figures can differ substantially.
The presence of wind energy, even when subsidized, can reduce costs for consumers (€5 billion/yr in Germany) by reducing the marginal price, by minimizing the use of expensive peaking power plants.
The cost has decreased as wind turbine technology has improved. There are now longer and lighter wind turbine blades, improvements in turbine performance, and increased power generation efficiency. Also, wind project capital expenditure costs and maintenance costs have continued to decline.
In 2021, a Lazard study of unsubsidized electricity said that wind power levelized cost of electricity continues to fall but more slowly than before. The study estimated new wind-generated electricity cost from $26 to $50/MWh, compared to new gas power from $45 to $74/MWh. The median cost of fully deprecated existing coal power was $42/MWh, nuclear $29/MWh and gas $24/MWh. The study estimated offshore wind at around $83/MWh. Compound annual growth rate was 4% per year from 2016 to 2021, compared to 10% per year from 2009 to 2021.
The value of wind power
While the levelised costs of wind power may have reached that of traditional combustion based power technologies, the market value of the generated power is also lower due to the merit order effect, which implies that electricity market prices are lower in hours with substantial generation of variable renewable energy due to the low marginal costs of this technology. The effect has been identified in several European markets. For wind power plants exposed to electricity market pricing in markets with high penetration of variable renewable energy sources, profitability can be challenged.
Incentives and community benefits
Turbine prices have fallen significantly in recent years due to tougher competitive conditions such as the increased use of energy auctions, and the elimination of subsidies in many markets. As of 2021, subsidies are still often given to offshore wind. But they are generally no longer necessary for onshore wind in countries with even a very low carbon price such as China, provided there are no competing fossil fuel subsidies.
Secondary market forces provide incentives for businesses to use wind-generated power, even if there is a premium price for the electricity. For example, socially responsible manufacturers pay utility companies a premium that goes to subsidize and build new wind power infrastructure. Companies use wind-generated power, and in return, they can claim that they are undertaking strong "green" efforts. Wind projects provide local taxes, or payments in place of taxes and strengthen the economy of rural communities by providing income to farmers with wind turbines on their land.
The wind energy sector can also produce jobs during the construction and operating phase. Jobs include the manufacturing of wind turbines and the construction process, which includes transporting, installing, and then maintaining the turbines. An estimated 1.25 million people were employed in wind power in 2020.
Small-scale wind power
Small-scale wind power is the name given to wind generation systems with the capacity to produce up to 50 kW of electrical power. Isolated communities, that may otherwise rely on diesel generators, may use wind turbines as an alternative. Individuals may purchase these systems to reduce or eliminate their dependence on grid electric power for economic reasons, or to reduce their carbon footprint. Wind turbines have been used for household electric power generation in conjunction with battery storage over many decades in remote areas.
Examples of small-scale wind power projects in an urban setting can be found in New York City, where, since 2009, several building projects have capped their roofs with Gorlov-type helical wind turbines. Although the energy they generate is small compared to the buildings' overall consumption, they help to reinforce the building's 'green' credentials in ways that "showing people your high-tech boiler" cannot, with some of the projects also receiving the direct support of the New York State Energy Research and Development Authority.
Grid-connected domestic wind turbines may use grid energy storage, thus replacing purchased electric power with locally produced power when available. The surplus power produced by domestic microgenerators can, in some jurisdictions, be fed into the network and sold to the utility company, producing a retail credit for the microgenerators' owners to offset their energy costs.
Off-grid system users can either adapt to intermittent power or use batteries, photovoltaic, or diesel systems to supplement the wind turbine. Equipment such as parking meters, traffic warning signs, street lighting, or wireless Internet gateways may be powered by a small wind turbine, possibly combined with a photovoltaic system, that charges a small battery replacing the need for a connection to the power grid.
Airborne wind turbines, such as kites, can be used in places at risk of hurricanes, as they can be taken down in advance.
Impact on environment and landscape
The environmental impact of electricity generation from wind power is minor when compared to that of fossil fuel power. Wind turbines have some of the lowest life-cycle greenhouse-gas emissions of energy sources: far less greenhouse gas is emitted than for the average unit of electricity, so wind power helps limit climate change. Use of engineered wood may allow carbon negative wind power. Wind power consumes no fuel, and emits no local air pollution, unlike fossil fuel power sources.
Onshore wind farms can have a significant visual impact. Due to a very low surface power density and spacing requirements, wind farms typically need to be spread over more land than other power stations. Their network of turbines, access roads, transmission lines, and substations can result in "energy sprawl"; although land between the turbines and roads can still be used for agriculture. Some wind farms are opposed for potentially spoiling protected scenic areas, archaeological landscapes and heritage sites. A report by the Mountaineering Council of Scotland concluded that wind farms harmed tourism in areas known for natural landscapes and panoramic views.
Habitat loss and fragmentation are the greatest potential impacts on wildlife of onshore wind farms, but the worldwide ecological impact is minimal. Thousands of birds and bats, including rare species, have been killed by wind turbine blades, though wind turbines are responsible for far fewer bird deaths than fossil-fueled power stations when climate change effects are included. Not including these effects, modern wind turbines kill about 0.273 birds per GWh in comparison with 0.200 by coal power plants. The effects of wind turbines on birds can be mitigated with proper wildlife monitoring.
Many wind turbine blades are made of fiberglass, and have a lifetime of 20 years. Blades are hollow: some blades are crushed to reduce their volume and then landfilled. However, as they can take a lot of weight they can be made into long lasting small bridges for walkers or cyclists. Blade end-of-life is complicated, and blades manufactured in the 2020s are more likely to be designed to be completely recyclable.
Wind turbines also generate noise. At a distance of , this may be around 45 dB, which is slightly louder than a refrigerator. At , they become inaudible. There are anecdotal reports of negative health effects on people who live very close to wind turbines. Peer-reviewed research has generally not supported these claims.
Politics
Central government
Although wind turbines with fixed bases are a mature technology and new installations are generally no longer subsidized, floating wind turbines are a relatively new technology so some governments subsidize them, for example to use deeper waters.
Fossil fuel subsidies by some governments are slowing the growth of renewables.
Permitting of wind farms can take years and some governments are trying to speed up – the wind industry says this will help limit climate change and increase energy security – sometimes groups such as fishers resist this but governments say that rules protecting biodiversity will still be followed.
Public opinion
Surveys of public attitudes across Europe and in many other countries show strong public support for wind power. Bakker et al. (2012) found in their study that residents who did not want turbines built near them suffered significantly more stress than those who "benefited economically from wind turbines".
Although wind power is a popular form of energy generation, onshore or near offshore wind farms are sometimes opposed for their impact on the landscape (especially scenic areas, heritage areas and archaeological landscapes), as well as noise, and impact on tourism.
In other cases, there is direct community ownership of wind farms. The hundreds of thousands of people who have become involved in Germany's small and medium-sized wind farms demonstrate such support there.
A 2010 Harris Poll found strong support for wind power in Germany, other European countries, and the United States.
Public support in the United States has decreased from 75% in 2020 to 62% in 2021, with the Democratic Party supporting the use of wind energy twice as much as the Republican Party. President Biden has signed an executive order to begin building large scale wind farms.
In China, Shen et al. (2019) found that Chinese city-dwellers may be resistant to building wind turbines in urban areas, with a surprisingly high proportion of people citing an unfounded fear of radiation as driving their concerns. Also, the study finds that like their counterparts in OECD countries, urban Chinese respondents are sensitive to direct costs and wildlife externalities. Distributing relevant information about turbines to the public may alleviate resistance.
Community
Many wind power companies work with local communities to reduce environmental and other concerns associated with particular wind farms.
In other cases there is direct community ownership of wind farm projects. Appropriate government consultation, planning and approval procedures also help to minimize environmental risks.
Some may still object to wind farms but many say their concerns should be weighed against the need to address the threats posed by air pollution, climate change and the opinions of the broader community.
In the US, wind power projects are reported to boost local tax bases, helping to pay for schools, roads, and hospitals, and to revitalize the economies of rural communities by providing steady income to farmers and other landowners.
In the UK, both the National Trust and the Campaign to Protect Rural England have expressed concerns about the effects on the rural landscape caused by inappropriately sited wind turbines and wind farms.
Some wind farms have become tourist attractions. The Whitelee Wind Farm Visitor Centre has an exhibition room, a learning hub, a café with a viewing deck and also a shop. It is run by the Glasgow Science Centre.
In Denmark, a loss-of-value scheme gives people the right to claim compensation for loss of value of their property if it is caused by proximity to a wind turbine. The loss must be at least 1% of the property's value.
Despite this general support for the concept of wind power in the public at large, local opposition often exists and has delayed or aborted a number of projects.
As well as concerns about the landscape, there are concerns that some installations can produce excessive sound and vibration levels leading to a decrease in property values. A study of 50,000 home sales near wind turbines found no statistical evidence that prices were affected.
While aesthetic issues are subjective and some find wind farms pleasant and optimistic, or symbols of energy independence and local prosperity, protest groups are often formed to attempt to block some wind power stations for various reasons.
Some opposition to wind farms is dismissed as NIMBYism, but research carried out in 2009 found that there is little evidence to support the belief that residents only object to wind farms because of a "Not in my Back Yard" attitude.
Geopolitics
Wind cannot be cut off unlike oil and gas so can contribute to energy security.
Turbine design
Wind turbines are devices that convert the wind's kinetic energy into electrical power. The result of over a millennium of windmill development and modern engineering, today's wind turbines are manufactured in a wide range of horizontal axis and vertical axis types. The smallest turbines are used for applications such as battery charging for auxiliary power. Slightly larger turbines can be used for making small contributions to a domestic power supply while selling unused power back to the utility supplier via the electrical grid. Arrays of large turbines, known as wind farms, have become an increasingly important source of renewable energy and are used in many countries as part of a strategy to reduce their reliance on fossil fuels.
Wind turbine design is the process of defining the form and specifications of a wind turbine to extract energy from the wind.
A wind turbine installation consists of the necessary systems needed to capture the wind's energy, point the turbine into the wind, convert mechanical rotation into electrical power, and other systems to start, stop, and control the turbine.
In 1919, the German physicist Albert Betz showed that for a hypothetical ideal wind-energy extraction machine, the fundamental laws of conservation of mass and energy allowed no more than 16/27 (59%) of the kinetic energy of the wind to be captured. This Betz limit can be approached in modern turbine designs, which may reach 70 to 80% of the theoretical Betz limit.
The aerodynamics of a wind turbine are not straightforward. The airflow at the blades is not the same as the airflow far away from the turbine. The very nature of how energy is extracted from the air also causes air to be deflected by the turbine. This affects the objects or other turbines downstream, which is known as "wake effect". Also, the aerodynamics of a wind turbine at the rotor surface exhibit phenomena that are rarely seen in other aerodynamic fields. The shape and dimensions of the blades of the wind turbine are determined by the aerodynamic performance required to efficiently extract energy from the wind, and by the strength required to resist the forces on the blade.
In addition to the aerodynamic design of the blades, the design of a complete wind power system must also address the design of the installation's rotor hub, nacelle, tower structure, generator, controls, and foundation.
History
Wind power has been used as long as humans have put sails into the wind. Wind-powered machines used to grind grain and pump water, the windmill and wind pump, were developed in what is now Iran, Afghanistan, and Pakistan by the 9th century. Wind power was widely available and not confined to the banks of fast-flowing streams, or later, requiring sources of fuel. Wind-powered pumps drained the polders of the Netherlands, and in arid regions such as the American mid-west or the Australian outback, wind pumps provided water for livestock and steam engines.
The first windmill used for the production of electric power was built in Scotland in July 1887 by Prof James Blyth of Anderson's College, Glasgow (the precursor of Strathclyde University). Blyth's high cloth-sailed wind turbine was installed in the garden of his holiday cottage at Marykirk in Kincardineshire, and was used to charge accumulators developed by the Frenchman Camille Alphonse Faure, to power the lighting in the cottage, thus making it the first house in the world to have its electric power supplied by wind power. Blyth offered the surplus electric power to the people of Marykirk for lighting the main street, however, they turned down the offer as they thought electric power was "the work of the devil." Although he later built a wind turbine to supply emergency power to the local Lunatic Asylum, Infirmary, and Dispensary of Montrose, the invention never really caught on as the technology was not considered to be economically viable.
Across the Atlantic, in Cleveland, Ohio, a larger and heavily engineered machine was designed and constructed in the winter of 1887–1888 by Charles F. Brush. This was built by his engineering company at his home and operated from 1886 until 1900. The Brush wind turbine had a rotor in diameter and was mounted on an tower. Although large by today's standards, the machine was only rated at 12 kW. The connected dynamo was used either to charge a bank of batteries or to operate up to 100 incandescent light bulbs, three arc lamps, and various motors in Brush's laboratory.
With the development of electric power, wind power found new applications in lighting buildings remote from centrally generated power. Throughout the 20th century parallel paths developed small wind stations suitable for farms or residences.
From 1932 many isolated properties in Australia ran their lighting and electric fans from batteries, charged by a "Freelite" wind-driven generator, producing 100watts of electrical power from as little wind speed as .
The 1973 oil crisis triggered the investigation in Denmark and the United States that led to larger utility-scale wind generators that could be connected to electric power grids for remote use of power. By 2008, the U.S. installed capacity had reached 25.4 gigawatts, and by 2012 the installed capacity was 60 gigawatts. Today, wind-powered generators operate in every size range between tiny stations for battery charging at isolated residences, up to gigawatt-sized offshore wind farms that provide electric power to national electrical networks. The European Union is working to augment these prospects.
In 2023, the global wind power sector experienced significant growth, with 116.6 gigawatts (GW) of new capacity added to the power grid, representing a 50% increase over the amount added in 2022. This surge in capacity brought the total installed wind power capacity worldwide to 1,021 GW by the end of the year, marking a growth of 13% compared to the previous year.
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277295 | https://en.wikipedia.org/wiki/Magnetite | Magnetite | Magnetite is a mineral and one of the main iron ores, with the chemical formula . It is one of the oxides of iron, and is ferrimagnetic; it is attracted to a magnet and can be magnetized to become a permanent magnet itself. With the exception of extremely rare native iron deposits, it is the most magnetic of all the naturally occurring minerals on Earth. Naturally magnetized pieces of magnetite, called lodestone, will attract small pieces of iron, which is how ancient peoples first discovered the property of magnetism.
Magnetite is black or brownish-black with a metallic luster, has a Mohs hardness of 5–6 and leaves a black streak. Small grains of magnetite are very common in igneous and metamorphic rocks.
The chemical IUPAC name is iron(II,III) oxide and the common chemical name is ferrous-ferric oxide.
Properties
In addition to igneous rocks, magnetite also occurs in sedimentary rocks, including banded iron formations and in lake and marine sediments as both detrital grains and as magnetofossils. Magnetite nanoparticles are also thought to form in soils, where they probably oxidize rapidly to maghemite.
Crystal structure
The chemical composition of magnetite is Fe2+(Fe3+)2(O2-)4. This indicates that magnetite contains both ferrous (divalent) and ferric (trivalent) iron, suggesting crystallization in an environment containing intermediate levels of oxygen. The main details of its structure were established in 1915. It was one of the first crystal structures to be obtained using X-ray diffraction. The structure is inverse spinel, with O2- ions forming a face-centered cubic lattice and iron cations occupying interstitial sites. Half of the Fe3+ cations occupy tetrahedral sites while the other half, along with Fe2+ cations, occupy octahedral sites. The unit cell consists of thirty-twoO2- ions and unit cell length is a = 0.839 nm.
As a member of the inverse spinel group, magnetite can form solid solutions with similarly structured minerals, including ulvospinel () and magnesioferrite ().
Titanomagnetite, also known as titaniferous magnetite, is a solid solution between magnetite and ulvospinel that crystallizes in many mafic igneous rocks. Titanomagnetite may undergo oxy-exsolution during cooling, resulting in ingrowths of magnetite and ilmenite.
Crystal morphology and size
Natural and synthetic magnetite occurs most commonly as octahedral crystals bounded by {111} planes and as rhombic-dodecahedra. Twinning occurs on the {111} plane.
Hydrothermal synthesis usually produces single octahedral crystals which can be as large as across. In the presence of mineralizers such as 0.1M HI or 2M NH4Cl and at 0.207MPa at 416–800 °C, magnetite grew as crystals whose shapes were a combination of rhombic-dodechahedra forms. The crystals were more rounded than usual. The appearance of higher forms was considered as a result from a decrease in the surface energies caused by the lower surface to volume ratio in the rounded crystals.
Reactions
Magnetite has been important in understanding the conditions under which rocks form. Magnetite reacts with oxygen to produce hematite, and the mineral pair forms a buffer that can control how oxidizing its environment is (the oxygen fugacity). This buffer is known as the hematite-magnetite or HM buffer. At lower oxygen levels, magnetite can form a buffer with quartz and fayalite known as the QFM buffer. At still lower oxygen levels, magnetite forms a buffer with wüstite known as the MW buffer. The QFM and MW buffers have been used extensively in laboratory experiments on rock chemistry. The QFM buffer, in particular, produces an oxygen fugacity close to that of most igneous rocks.
Commonly, igneous rocks contain solid solutions of both titanomagnetite and hemoilmenite or titanohematite. Compositions of the mineral pairs are used to calculate oxygen fugacity: a range of oxidizing conditions are found in magmas and the oxidation state helps to determine how the magmas might evolve by fractional crystallization. Magnetite also is produced from peridotites and dunites by serpentinization.
Magnetic properties
Lodestones were used as an early form of magnetic compass. Magnetite has been a critical tool in paleomagnetism, a science important in understanding plate tectonics and as historic data for magnetohydrodynamics and other scientific fields.
The relationships between magnetite and other iron oxide minerals such as ilmenite, hematite, and ulvospinel have been much studied; the reactions between these minerals and oxygen influence how and when magnetite preserves a record of the Earth's magnetic field.
At low temperatures, magnetite undergoes a crystal structure phase transition from a monoclinic structure to a cubic structure known as the Verwey transition. Optical studies show that this metal to insulator transition is sharp and occurs around 120K. The Verwey transition is dependent on grain size, domain state, pressure, and the iron-oxygen stoichiometry. An isotropic point also occurs near the Verwey transition around 130K, at which point the sign of the magnetocrystalline anisotropy constant changes from positive to negative. The Curie temperature of magnetite is .
If magnetite is in a large enough quantity it can be found in aeromagnetic surveys using a magnetometer which measures magnetic intensities.
Melting point
Solid magnetite particles melt at about .
Distribution of deposits
Magnetite is sometimes found in large quantities in beach sand. Such black sands (mineral sands or iron sands) are found in various places, such as Lung Kwu Tan in Hong Kong; California, United States; and the west coast of the North Island of New Zealand. The magnetite, eroded from rocks, is carried to the beach by rivers and concentrated by wave action and currents. Huge deposits have been found in banded iron formations. These sedimentary rocks have been used to infer changes in the oxygen content of the atmosphere of the Earth.
Large deposits of magnetite are also found in the Atacama region of Chile (Chilean Iron Belt); the Valentines region of Uruguay; Kiruna, Sweden; the Tallawang region of New South Wales; and in the Adirondack Mountains of New York in the United States. Kediet ej Jill, the highest mountain of Mauritania, is made entirely of the mineral. In the municipalities of Molinaseca, Albares, and Rabanal del Camino, in the province of León (Spain), there is a magnetite deposit in Ordovician terrain, considered one of the largest in Europe. It was exploited between 1955 and 1982. Deposits are also found in Norway, Romania, and Ukraine. Magnetite-rich sand dunes are found in southern Peru. In 2005, an exploration company, Cardero Resources, discovered a vast deposit of magnetite-bearing sand dunes in Peru. The dune field covers 250 square kilometers (100 sq mi), with the highest dune at over 2,000 meters (6,560 ft) above the desert floor. The sand contains 10% magnetite.
In large enough quantities magnetite can affect compass navigation. In Tasmania there are many areas with highly magnetized rocks that can greatly influence compasses. Extra steps and repeated observations are required when using a compass in Tasmania to keep navigation problems to the minimum.
Magnetite crystals with a cubic habit are rare but have been found at Balmat, St. Lawrence County, New York, and at Långban, Sweden. This habit may be a result of crystallization in the presence of cations such as zinc.
Magnetite can also be found in fossils due to biomineralization and are referred to as magnetofossils. There are also instances of magnetite with origins in space coming from meteorites.
Biological occurrences
Biomagnetism is usually related to the presence of biogenic crystals of magnetite, which occur widely in organisms. These organisms range from magnetotactic bacteria (e.g., Magnetospirillum magnetotacticum) to animals, including humans, where magnetite crystals (and other magnetically sensitive compounds) are found in different organs, depending on the species. Biomagnetites account for the effects of weak magnetic fields on biological systems. There is also a chemical basis for cellular sensitivity to electric and magnetic fields (galvanotaxis).
Pure magnetite particles are biomineralized in magnetosomes, which are produced by several species of magnetotactic bacteria. Magnetosomes consist of long chains of oriented magnetite particle that are used by bacteria for navigation. After the death of these bacteria, the magnetite particles in magnetosomes may be preserved in sediments as magnetofossils. Some types of anaerobic bacteria that are not magnetotactic can also create magnetite in oxygen free sediments by reducing amorphic ferric oxide to magnetite.
Several species of birds are known to incorporate magnetite crystals in the upper beak for magnetoreception, which (in conjunction with cryptochromes in the retina) gives them the ability to sense the direction, polarity, and magnitude of the ambient magnetic field.
Chitons, a type of mollusk, have a tongue-like structure known as a radula, covered with magnetite-coated teeth, or denticles. The hardness of the magnetite helps in breaking down food.
Biological magnetite may store information about the magnetic fields the organism was exposed to, potentially allowing scientists to learn about the migration of the organism or about changes in the Earth's magnetic field over time.
Human brain
Living organisms can produce magnetite. In humans, magnetite can be found in various parts of the brain including the frontal, parietal, occipital, and temporal lobes, brainstem, cerebellum and basal ganglia. Iron can be found in three forms in the brain – magnetite, hemoglobin (blood) and ferritin (protein), and areas of the brain related to motor function generally contain more iron. Magnetite can be found in the hippocampus. The hippocampus is associated with information processing, specifically learning and memory. However, magnetite can have toxic effects due to its charge or magnetic nature and its involvement in oxidative stress or the production of free radicals. Research suggests that beta-amyloid plaques and tau proteins associated with neurodegenerative disease frequently occur after oxidative stress and the build-up of iron.
Some researchers also suggest that humans possess a magnetic sense, proposing that this could allow certain people to use magnetoreception for navigation. The role of magnetite in the brain is still not well understood, and there has been a general lag in applying more modern, interdisciplinary techniques to the study of biomagnetism.
Electron microscope scans of human brain-tissue samples are able to differentiate between magnetite produced by the body's own cells and magnetite absorbed from airborne pollution, the natural forms being jagged and crystalline, while magnetite pollution occurs as rounded nanoparticles. Potentially a human health hazard, airborne magnetite is a result of pollution (specifically combustion). These nanoparticles can travel to the brain via the olfactory nerve, increasing the concentration of magnetite in the brain. In some brain samples, the nanoparticle pollution outnumbers the natural particles by as much as 100:1, and such pollution-borne magnetite particles may be linked to abnormal neural deterioration. In one study, the characteristic nanoparticles were found in the brains of 37 people: 29 of these, aged 3 to 85, had lived and died in Mexico City, a significant air pollution hotspot. Some of the further eight, aged 62 to 92, from Manchester, England, had died with varying severities of neurodegenerative diseases. Such particles could conceivably contribute to diseases like Alzheimer's disease. Though a causal link has not yet been established, laboratory studies suggest that iron oxides such as magnetite are a component of protein plaques in the brain. Such plaques have been linked to Alzheimer's disease.
Increased iron levels, specifically magnetic iron, have been found in portions of the brain in Alzheimer's patients. Monitoring changes in iron concentrations may make it possible to detect the loss of neurons and the development of neurodegenerative diseases prior to the onset of symptoms due to the relationship between magnetite and ferritin. In tissue, magnetite and ferritin can produce small magnetic fields which will interact with magnetic resonance imaging (MRI) creating contrast. Huntington patients have not shown increased magnetite levels; however, high levels have been found in study mice.
Applications
Due to its high iron content, magnetite has long been a major iron ore. It is reduced in blast furnaces to pig iron or sponge iron for conversion to steel.
Magnetic recording
Audio recording using magnetic acetate tape was developed in the 1930s. The German magnetophon first utilized magnetite powder that BASF coated onto cellulose acetate before soon switching to gamma ferric oxide for its superior morphology. Following World War II, 3M Company continued work on the German design. In 1946, the 3M researchers found they could also improve their own magnetite-based paper tape, which utilized powders of cubic crystals, by replacing the magnetite with needle-shaped particles of gamma ferric oxide (γ-Fe2O3).
Catalysis
Approximately 2–3% of the world's energy budget is allocated to the Haber Process for nitrogen fixation, which relies on magnetite-derived catalysts. The industrial catalyst is obtained from finely ground iron powder, which is usually obtained by reduction of high-purity magnetite. The pulverized iron metal is burnt (oxidized) to give magnetite or wüstite of a defined particle size. The magnetite (or wüstite) particles are then partially reduced, removing some of the oxygen in the process. The resulting catalyst particles consist of a core of magnetite, encased in a shell of wüstite, which in turn is surrounded by an outer shell of iron metal. The catalyst maintains most of its bulk volume during the reduction, resulting in a highly porous high-surface-area material, which enhances its effectiveness as a catalyst.
Magnetite nanoparticles
Magnetite micro- and nanoparticles are used in a variety of applications, from biomedical to environmental. One use is in water purification: in high gradient magnetic separation, magnetite nanoparticles introduced into contaminated water will bind to the suspended particles (solids, bacteria, or plankton, for example) and settle to the bottom of the fluid, allowing the contaminants to be removed and the magnetite particles to be recycled and reused. This method works with radioactive and carcinogenic particles as well, making it an important cleanup tool in the case of heavy metals introduced into water systems.
Another application of magnetic nanoparticles is in the creation of ferrofluids. These are used in several ways. Ferrofluids can be used for targeted drug delivery in the human body. The magnetization of the particles bound with drug molecules allows "magnetic dragging" of the solution to the desired area of the body. This would allow the treatment of only a small area of the body, rather than the body as a whole, and could be highly useful in cancer treatment, among other things. Ferrofluids are also used in magnetic resonance imaging (MRI) technology.
Coal mining industry
For the separation of coal from waste, dense medium baths were used. This technique employed the difference in densities between coal (1.3–1.4 tonnes per m3) and shales (2.2–2.4 tonnes per m3). In a medium with intermediate density (water with magnetite), stones sank and coal floated.
Magnetene
Magnetene is a two-dimensional flat sheet of magnetite noted for its ultra-low-friction properties.
Gallery
| Physical sciences | Minerals | Earth science |
277353 | https://en.wikipedia.org/wiki/Pier | Pier | A pier is a raised structure that rises above a body of water and usually juts out from its shore, typically supported by piles or pillars, and provides above-water access to offshore areas. Frequent pier uses include fishing, boat docking and access for both passengers and cargo, and oceanside recreation. Bridges, buildings, and walkways may all be supported by architectural piers. Their open structure allows tides and currents to flow relatively unhindered, whereas the more solid foundations of a quay or the closely spaced piles of a wharf can act as a breakwater, and are consequently more liable to silting. Piers can range in size and complexity from a simple lightweight wooden structure to major structures extended over . In American English, a pier may be synonymous with a dock.
Piers have been built for several purposes, and because these different purposes have distinct regional variances, the term pier tends to have different nuances of meaning in different parts of the world. Thus in North America and Australia, where many ports were, until recently, built on the multiple pier model, the term tends to imply a current or former cargo-handling facility. In contrast, in Europe, where ports more often use basins and river-side quays than piers, the term is principally associated with the image of a Victorian cast iron pleasure pier which emerged in Great Britain during the early 19th century. However, the earliest piers pre-date the Victorian age.
Types
Piers can be categorized into different groupings according to the principal purpose. However, there is considerable overlap between these categories. For example, pleasure piers often also allow for the docking of pleasure steamers and other similar craft, while working piers have often been converted to leisure use after being rendered obsolete by advanced developments in cargo-handling technology. Many piers are floating piers, to ensure that the piers raise and lower with the tide along with the boats tied to them. This prevents a situation where lines become overly taut or loose by rising or lowering tides. An overly taut or loose tie-line can damage boats by pulling them out of the water or allowing them so much leeway that they bang forcefully against the sides of the pier.
Working piers
Working piers were built for the handling of passengers and cargo onto and off ships or (as at Wigan Pier) canal boats. Working piers themselves fall into two different groups. Longer individual piers are often found at ports with large tidal ranges, with the pier stretching far enough off shore to reach deep water at low tide. Such piers provided an economical alternative to impounded docks where cargo volumes were low, or where specialist bulk cargo was handled, such as at coal piers. The other form of working pier, often called the finger pier, was built at ports with smaller tidal ranges. Here the principal advantage was to give a greater available quay length for ships to berth against compared to a linear littoral quayside, and such piers are usually much shorter. Typically each pier would carry a single transit shed the length of the pier, with ships berthing bow or stern in to the shore. Some major ports consisted of large numbers of such piers lining the foreshore, classic examples being the Hudson River frontage of New York, or the Embarcadero in San Francisco.
The advent of container shipping, with its need for large container handling spaces adjacent to the shipping berths, has made working piers obsolete for the handling of general cargo, although some still survive for the handling of passenger ships or bulk cargos. One example, is in use in Progreso, Yucatán, where a pier extends more than 4 miles into the Gulf of Mexico, making it the longest pier in the world. The Progreso Pier supplies much of the peninsula with transportation for the fishing and cargo industries and serves as a port for large cruise ships in the area. Many other working piers have been demolished, or remain derelict, but some have been recycled as pleasure piers. The best known example of this is Pier 39 in San Francisco.
At Southport and the Tweed River on the Gold Coast in Australia, there are piers that support equipment for a sand bypassing system that maintains the health of sandy beaches and navigation channels.
Pleasure piers
Pleasure piers were first built in Britain during the early 19th century. The earliest structures were Ryde Pier, built in 1813/4, Trinity Chain Pier near Leith, built in 1821, Brighton Chain Pier, built in 1823. and Margate Jetty 1823/24 originally a timber built pier.
Only the oldest of these piers still remains. At that time, the introduction of steamships and railways for the first time permitted mass tourism to dedicated seaside resorts. The large tidal ranges at many such resorts meant that passengers arriving by pleasure steamer could use a pier to disembark safely. Also, for much of the day, the sea was not visible from the shore and the pleasure pier permitted holidaymakers to promenade over and alongside the sea at all times. The world's longest pleasure pier is at Southend-on-Sea, Essex, and extends 1.3 miles (2.1 km) into the Thames Estuary. The longest pier on the West Coast of the US is the Santa Cruz Wharf, with a length of .
Providing a walkway out to sea, pleasure piers often include amusements and theatres as part of their attractions. Such a pier may be unroofed, closed, or partly open and partly closed. Sometimes a pier has two decks. Galveston Island Historic Pleasure Pier in Galveston, Texas has a roller coaster, 15 rides, carnival games and souvenir shops.
Early pleasure piers were of complete timber construction, as was with Margate which opened in 1824. The first iron and timber built pleasure pier Margate Jetty, opened in 1855. Margate pier was wrecked by a storm in January 1978 and not repaired. The longest iron pleasure pier still remaining is the one at Southend. First opened as a wooden pier in 1829, it was reconstructed in iron and completed in 1889. In a 2006 UK poll, the public voted the seaside pier onto the list of icons of England.
Fishing piers
Many piers are built for the purpose of providing boatless anglers access to fishing grounds that are otherwise inaccessible. Many "Free Piers" are available in larger harbors which differ from private piers. Free Piers are often primarily used for fishing. Fishing from a pier presents a set of different circumstances to fishing from the shore or beach, as you do not need to cast out into the deeper water. This being the case there are specific fishing rigs that have been created specifically for pier fishing which allow for the direct access to deeper water.
Piers of the world
Belgium
In Blankenberge a first pleasure pier was built in 1894. After its destruction in the World War I, a new pier was built in 1933. It remained till the present day, but was partially transformed and modernized in 1999–2004.
In Nieuwpoort, Belgium there is a pleasure pier on both sides of the river IJzer.
Netherlands
Scheveningen, the coastal resort town of The Hague, boasts the largest pier in the Netherlands, completed in 1961. A crane, built on top of the pier's panorama tower, provides the opportunity to make a high bungee jump over the North Sea waves. The present pier is a successor of an earlier pier, which was completed in 1901 but in 1943 destroyed by the German occupation forces.
United Kingdom
England and Wales
The first recorded pier in England was Ryde Pier, opened in 1814 on the Isle of Wight, as a landing stage to allow ferries to and from the mainland to berth. It is still used for this purpose today. It also had a leisure function in the past, with the pier head once containing a pavilion, and there are still refreshment facilities today. The oldest cast iron pier in the world is Town Pier, Gravesend, in Kent, which opened in 1834. However, it is not recognised by the National Piers Society as being a seaside pier.
Following the building of the world's first seaside pier at Ryde, the pier became fashionable at seaside resorts in England and Wales during the Victorian era, peaking in the 1860s with 22 being built in that decade. A symbol of the typical British seaside holiday, by 1914, more than 100 pleasure piers were located around the UK coast. Regarded as being among the finest Victorian architecture, there are still a significant number of seaside piers of architectural merit still standing, although some have been lost, including Margate, two at Brighton in East Sussex, one at New Brighton in the Wirral and three at Blackpool in Lancashire. Two piers, Brighton's now derelict West Pier and Clevedon Pier, were Grade 1 listed. The Birnbeck Pier in Weston-super-Mare is the only pier in the world linked to an island. The National Piers Society gives a figure of 55 surviving seaside piers in England and Wales. In 2017, Brighton Palace Pier was said to be the most visited tourist attraction outside London, with over 4.5 million visitors the previous year.
| Technology | Coastal infrastructure | null |
277379 | https://en.wikipedia.org/wiki/Margin%20of%20error | Margin of error | The margin of error is a statistic expressing the amount of random sampling error in the results of a survey. The larger the margin of error, the less confidence one should have that a poll result would reflect the result of a census of the entire population. The margin of error will be positive whenever a population is incompletely sampled and the outcome measure has positive variance, which is to say, whenever the measure varies.
The term margin of error is often used in non-survey contexts to indicate observational error in reporting measured quantities.
Concept
Consider a simple yes/no poll as a sample of respondents drawn from a population reporting the percentage of yes responses. We would like to know how close is to the true result of a survey of the entire population , without having to conduct one. If, hypothetically, we were to conduct a poll over subsequent samples of respondents (newly drawn from ), we would expect those subsequent results to be normally distributed about , the true but unknown percentage of the population. The margin of error describes the distance within which a specified percentage of these results is expected to vary from .
Going by the Central limit theorem, the margin of error helps to explain how the distribution of sample means (or percentage of yes, in this case) will approximate a normal distribution as sample size increases. If this applies, it would speak about the sampling being unbiased, but not about the inherent distribution of the data.
According to the 68-95-99.7 rule, we would expect that 95% of the results will fall within about two standard deviations () either side of the true mean . This interval is called the confidence interval, and the radius (half the interval) is called the margin of error, corresponding to a 95% confidence level.
Generally, at a confidence level , a sample sized of a population having expected standard deviation has a margin of error
where denotes the quantile (also, commonly, a z-score), and is the standard error.
Standard deviation and standard error
We would expect the average of normally distributed values to have a standard deviation which somehow varies with . The smaller , the wider the margin. This is called the standard error .
For the single result from our survey, we assume that , and that all subsequent results together would have a variance .
Note that corresponds to the variance of a Bernoulli distribution.
Maximum margin of error at different confidence levels
For a confidence level , there is a corresponding confidence interval about the mean , that is, the interval within which values of should fall with probability . Precise values of are given by the quantile function of the normal distribution (which the 68–95–99.7 rule approximates).
Note that is undefined for , that is, is undefined, as is .
Since at , we can arbitrarily set , calculate , , and to obtain the maximum margin of error for at a given confidence level and sample size , even before having actual results. With
Also, usefully, for any reported
Specific margins of error
If a poll has multiple percentage results (for example, a poll measuring a single multiple-choice preference), the result closest to 50% will have the highest margin of error. Typically, it is this number that is reported as the margin of error for the entire poll. Imagine poll reports as
(as in the figure above)
As a given percentage approaches the extremes of 0% or 100%, its margin of error approaches ±0%.
Comparing percentages
Imagine multiple-choice poll reports as . As described above, the margin of error reported for the poll would typically be , as is closest to 50%. The popular notion of statistical tie or statistical dead heat, however, concerns itself not with the accuracy of the individual results, but with that of the ranking of the results. Which is in first?
If, hypothetically, we were to conduct a poll over subsequent samples of respondents (newly drawn from ), and report the result , we could use the standard error of difference to understand how is expected to fall about . For this, we need to apply the sum of variances to obtain a new variance, ,
where is the covariance of and .
Thus (after simplifying),
Note that this assumes that is close to constant, that is, respondents choosing either A or B would almost never choose C (making and close to perfectly negatively correlated). With three or more choices in closer contention, choosing a correct formula for becomes more complicated.
Effect of finite population size
The formulae above for the margin of error assume that there is an infinitely large population and thus do not depend on the size of population , but only on the sample size . According to sampling theory, this assumption is reasonable when the sampling fraction is small. The margin of error for a particular sampling method is essentially the same regardless of whether the population of interest is the size of a school, city, state, or country, as long as the sampling fraction is small.
In cases where the sampling fraction is larger (in practice, greater than 5%), analysts might adjust the margin of error using a finite population correction to account for the added precision gained by sampling a much larger percentage of the population. FPC can be calculated using the formula
...and so, if poll were conducted over 24% of, say, an electorate of 300,000 voters,
Intuitively, for appropriately large ,
In the former case, is so small as to require no correction. In the latter case, the poll effectively becomes a census and sampling error becomes moot.
| Mathematics | Statistics | null |
277468 | https://en.wikipedia.org/wiki/Two-body%20problem | Two-body problem | In classical mechanics, the two-body problem is to calculate and predict the motion of two massive bodies that are orbiting each other in space. The problem assumes that the two bodies are point particles that interact only with one another; the only force affecting each object arises from the other one, and all other objects are ignored.
The most prominent example of the classical two-body problem is the gravitational case (see also Kepler problem), arising in astronomy for predicting the orbits (or escapes from orbit) of objects such as satellites, planets, and stars. A two-point-particle model of such a system nearly always describes its behavior well enough to provide useful insights and predictions.
A simpler "one body" model, the "central-force problem", treats one object as the immobile source of a force acting on the other. One then seeks to predict the motion of the single remaining mobile object. Such an approximation can give useful results when one object is much more massive than the other (as with a light planet orbiting a heavy star, where the star can be treated as essentially stationary).
However, the one-body approximation is usually unnecessary except as a stepping stone. For many forces, including gravitational ones, the general version of the two-body problem can be reduced to a pair of one-body problems, allowing it to be solved completely, and giving a solution simple enough to be used effectively.
By contrast, the three-body problem (and, more generally, the n-body problem for n ≥ 3) cannot be solved in terms of first integrals, except in special cases.
Results for prominent cases
Gravitation and other inverse-square examples
The two-body problem is interesting in astronomy because pairs of astronomical objects are often moving rapidly in arbitrary directions (so their motions become interesting), widely separated from one another (so they will not collide) and even more widely separated from other objects (so outside influences will be small enough to be ignored safely).
Under the force of gravity, each member of a pair of such objects will orbit their mutual center of mass in an elliptical pattern, unless they are moving fast enough to escape one another entirely, in which case their paths will diverge along other planar conic sections. If one object is very much heavier than the other, it will move far less than the other with reference to the shared center of mass. The mutual center of mass may even be inside the larger object.
For the derivation of the solutions to the problem, see Classical central-force problem or Kepler problem.
In principle, the same solutions apply to macroscopic problems involving objects interacting not only through gravity, but through any other attractive scalar force field obeying an inverse-square law, with electrostatic attraction being the obvious physical example. In practice, such problems rarely arise. Except perhaps in experimental apparatus or other specialized equipment, we rarely encounter electrostatically interacting objects which are moving fast enough, and in such a direction, as to avoid colliding, and/or which are isolated enough from their surroundings.
The dynamical system of a two-body system under the influence of torque turns out to be a Sturm-Liouville equation.
Inapplicability to atoms and subatomic particles
Although the two-body model treats the objects as point particles, classical mechanics only apply to systems of macroscopic scale. Most behavior of subatomic particles cannot be predicted under the classical assumptions underlying this article or using the mathematics here.
Electrons in an atom are sometimes described as "orbiting" its nucleus, following an early conjecture of Niels Bohr (this is the source of the term "orbital"). However, electrons don't actually orbit nuclei in any meaningful sense, and quantum mechanics are necessary for any useful understanding of the electron's real behavior. Solving the classical two-body problem for an electron orbiting an atomic nucleus is misleading and does not produce many useful insights.
Reduction to two independent, one-body problems
The complete two-body problem can be solved by re-formulating it as two one-body problems: a trivial one and one that involves solving for the motion of one particle in an external potential. Since many one-body problems can be solved exactly, the corresponding two-body problem can also be solved.
Let and be the vector positions of the two bodies, and m1 and m2 be their masses. The goal is to determine the trajectories and for all times t, given the initial positions and and the initial velocities and .
When applied to the two masses, Newton's second law states that
where F12 is the force on mass 1 due to its interactions with mass 2, and F21 is the force on mass 2 due to its interactions with mass 1. The two dots on top of the x position vectors denote their second derivative with respect to time, or their acceleration vectors.
Adding and subtracting these two equations decouples them into two one-body problems, which can be solved independently. Adding equations (1) and () results in an equation describing the center of mass (barycenter) motion. By contrast, subtracting equation (2) from equation (1) results in an equation that describes how the vector between the masses changes with time. The solutions of these independent one-body problems can be combined to obtain the solutions for the trajectories and .
Center of mass motion (1st one-body problem)
Let be the position of the center of mass (barycenter) of the system. Addition of the force equations (1) and (2) yields
where we have used Newton's third law and where
The resulting equation:
shows that the velocity of the center of mass is constant, from which follows that the total momentum is also constant (conservation of momentum). Hence, the position of the center of mass can be determined at all times from the initial positions and velocities.
Displacement vector motion (2nd one-body problem)
Dividing both force equations by the respective masses, subtracting the second equation from the first, and rearranging gives the equation
where we have again used Newton's third law and where is the displacement vector from mass 2 to mass 1, as defined above.
The force between the two objects, which originates in the two objects, should only be a function of their separation and not of their absolute positions and ; otherwise, there would not be translational symmetry, and the laws of physics would have to change from place to place. The subtracted equation can therefore be written:
where is the reduced mass
Solving the equation for is the key to the two-body problem. The solution depends on the specific force between the bodies, which is defined by . For the case where follows an inverse-square law, see the Kepler problem.
Once and have been determined, the original trajectories may be obtained
as may be verified by substituting the definitions of R and r into the right-hand sides of these two equations.
Two-body motion is planar
The motion of two bodies with respect to each other always lies in a plane (in the center of mass frame).
Proof: Defining the linear momentum and the angular momentum of the system, with respect to the center of mass, by the equations
where is the reduced mass and is the relative position (with these written taking the center of mass as the origin, and thus both parallel to ) the rate of change of the angular momentum equals the net torque
and using the property of the vector cross product that for any vectors and pointing in the same direction,
with .
Introducing the assumption (true of most physical forces, as they obey Newton's strong third law of motion) that the force between two particles acts along the line between their positions, it follows that and the angular momentum vector is constant (conserved). Therefore, the displacement vector and its velocity are always in the plane perpendicular to the constant vector .
Energy of the two-body system
If the force is conservative then the system has a potential energy , so the total energy can be written as
In the center of mass frame the kinetic energy is the lowest and the total energy becomes
The coordinates and can be expressed as
and in a similar way the energy E is related to the energies and that separately contain the kinetic energy of each body:
Central forces
For many physical problems, the force is a central force, i.e., it is of the form
where and is the corresponding unit vector. We now have:
where is negative in the case of an attractive force.
| Physical sciences | Classical mechanics | Physics |
277607 | https://en.wikipedia.org/wiki/Bull%20shark | Bull shark | The bull shark (Carcharhinus leucas), also known as the Zambezi shark (informally zambi) in Africa and Lake Nicaragua shark in Nicaragua, is a species of requiem shark commonly found worldwide in warm, shallow waters along coasts and in rivers. It is known for its aggressive nature, and presence mainly in warm, shallow brackish and freshwater systems including estuaries and (usually) lower reaches of rivers. Their aggressive nature has led to ongoing shark-culling efforts near beaches to protect beachgoers, which is one of the causes of bull shark populations continuing to decrease. Bull sharks are currently listed as vulnerable on the IUCN Red List.
Bull sharks are euryhaline and can thrive in both salt and fresh water. They are known to travel far up rivers, and have been known to travel up the Mississippi River as far as Alton, Illinois, about from the ocean, but few freshwater interactions with humans have been recorded. Larger-sized bull sharks are probably responsible for the majority of nearshore shark attacks, including many incidents of shark bites attributed to other species.
Unlike the river sharks of the genus Glyphis, bull sharks are not true freshwater sharks, despite their ability to survive in freshwater habitats.
This shark appears in the image of the 2000 colones bill from Costa Rica.
Etymology
The name "bull shark" comes from the shark's stocky shape, broad, flat snout, and aggressive, unpredictable behavior. In India, the bull shark may be confused with the Sundarbans or Ganges shark. In Africa, it is also commonly called the Zambezi River shark, or just "zambi".
Its wide range and diverse habitats result in many other local names, including Ganges River shark, Fitzroy Creek whaler, van Rooyen's shark, Lake Nicaragua shark, river shark, freshwater whaler, estuary whaler, Swan River whaler, cub shark, and shovelnose shark.
Evolution
Some of the bull shark's closest living relatives do not have the capabilities of osmoregulation. Its genus, Carcharhinus, also includes the sandbar shark, which is not capable of osmoregulation.
The bull shark shares numerous similarities with river sharks of the genus Glyphis, and other species in the genus Carcharhinus, but its phylogeny has not been cleared yet.
Anatomy and appearance
Bull sharks are large and stout, with females being larger than males. The bull shark can be up to in length at birth. Adult female bull sharks average long and typically weigh , whereas the slightly smaller adult male averages and . While a maximum size of is commonly reported, a single record exists of a female specimen of exactly . A long pregnant individual reached . The maximum weight can be over , ranking it among the largest of the requiem sharks. Bull sharks are wider and heavier than other requiem sharks of comparable length, and are grey on top and white below. The second dorsal fin is smaller than the first. The bull shark's caudal fin is longer and lower than that of the larger sharks, and it has a small snout, and lacks an interdorsal ridge.
Bull sharks have a bite force up to , weight for weight the highest among all investigated cartilaginous fishes.
Exceptional specimens
In early June 2012, off the coast of the Florida Keys near the western part of the Atlantic Ocean, a female believed to measure at least and was caught by members of the R.J. Dunlap Marine Conservation Program. In the Arabian Sea, off the coast of Fujairah in the United Arab Emirates, a pregnant shark weighing and measuring long was caught in February 2019, followed by another specimen weighing about and measuring about the same in length, in January 2020. Unconfirmed reports suggest that the very largest, exceptional specimens can possibly weigh up to .
Distribution and habitat
The bull shark is commonly found worldwide in coastal areas of warm oceans, in rivers and lakes, and occasionally salt and freshwater streams if they are deep enough. It is found to a depth of , but does not usually swim deeper than . In the Atlantic, it is found from Massachusetts to southern Brazil, and from Morocco to Angola.
Populations of bull sharks are also found in several major rivers, with more than 500 bull sharks thought to be living in the Brisbane River. One was reportedly seen swimming the flooded streets of Brisbane, Queensland, Australia, during the 2010–11 Queensland floods. Several were sighted in one of the main streets of Goodna, Queensland, shortly after the peak of the January 2011, floods. A large bull shark was caught in the canals of Scarborough, just north of Brisbane within Moreton Bay. Still greater numbers are in the canals of the Gold Coast, Queensland. In the warmer months of the year, bull sharks frequent Sydney Harbour. In the Pacific Ocean, it can be found from Baja California to Ecuador.
The bull shark has traveled up the Amazon River to Iquitos in Peru and north Bolivia. It also lives in freshwater Lake Nicaragua, in the Ganges and Brahmaputra Rivers of West Bengal, and Assam in Eastern India and adjoining Bangladesh. It can live in water with a high salt content as in St. Lucia Estuary in South Africa. Bull sharks have been recorded in the Tigris River since at least 1924 as far upriver as Baghdad, and has been rumored to also inhabit the Cahora Bassa lake upstream of the Zambezi. The species has a distinct preference for warm currents.
After Hurricane Katrina, many bull sharks were sighted in Lake Pontchartrain. In July 2023 some local fishermen in the Atchafalaya River have reported increasing numbers. Bull sharks have occasionally gone as far upstream in the Mississippi River as Alton, Illinois. Bull sharks have also been found in the Potomac River in St. Mary's County, Maryland. From 1996 to 2013, a golf course lake at Carbrook, Logan City, Queensland, Australia was the home to several bull sharks. They were trapped following a flood of the Logan and Albert Rivers in 1996, and resided in the lake until 2013, when they disappeared after another series of floods. The golf course capitalized on the novelty, changing their logo to feature the sharks and hosting a monthly tournament called the "Shark Lake Challenge".
Behavior
Freshwater tolerance
The bull shark is the best known of 43 species of elasmobranch, across 10 genera and four families, to have been reported in fresh and/or brackish water. Other species that enter rivers include the stingrays (Dasyatidae, Potamotrygonidae and others) and sawfish (Pristidae). Some skates (Rajidae), smooth dogfishes (Triakidae), and sandbar sharks (Carcharhinus plumbeus) regularly enter estuaries.
The bull shark is diadromous, meaning they can swim between salt and fresh water with ease, as they are euryhaline fish—able to quickly adapt to a wide range of salinities. Thus, the bull shark is one of the few cartilaginous fishes that have been reported in freshwater systems. Many of the euryhaline fish are bony fishes, such as salmon or tilapia, and are not closely related to bull sharks. Evolutionary assumptions can be made to help explain this sort of evolutionary disconnect, with one being that the bull shark experienced a population bottleneck during the last ice age. This bottleneck may have separated the bull shark from the rest of the Elasmobranchii subclass and favored the genes for an osmoregulatory system.
Elasmobranchs' ability to enter fresh water is limited because their blood is normally at least as salty (in terms of osmotic strength) as seawater through the accumulation of urea and trimethylamine oxide, but bull sharks living in fresh water show a significantly reduced concentration of urea within their blood. Despite this, the solute composition (i.e. osmolarity) of a bull shark in fresh water is still much higher than that of the external environment. This results in a large influx of water across the gills due to osmosis and loss of sodium and chloride from the shark's body. However, bull sharks in fresh water possess several organs with which to maintain appropriate salt and water balance; these are the rectal gland, kidneys, liver, and gills. All elasmobranchs have a rectal gland which functions in the excretion of excess salts accumulated as a consequence of living in seawater. Bull sharks in freshwater environments decrease the salt-excretory activity of the rectal gland, thereby conserving sodium and chloride. The kidneys produce large amounts of dilute urine, but also play an important role in the active reabsorption of solutes into the blood. The gills of bull sharks are likely to be involved in the uptake of sodium and chloride from the surrounding fresh water, whereas urea is produced in the liver as required with changes in environmental salinity. Recent work also shows that the differences in density of fresh water to that of marine waters result in significantly greater negative buoyancies in sharks occupying fresh water, resulting in increasing costs of living in fresh water. Bull sharks caught in freshwater have subsequently been shown to have lower liver densities than sharks living in marine waters. This may reduce the added cost of greater negative buoyancy.
Bull sharks are able to regulate themselves to live in either fresh or salt water. It can live in fresh water for its entire life, but this does not happen, mostly due to the reproductive needs of the shark. Young bull sharks leave the brackish water in which they are born and move out into the sea to breed. While is theoretically possible for bull sharks to live purely in fresh water, experiments conducted on bull sharks found that they died within four years. The stomach was opened and all that was found were two small, unidentifiable fishes. The cause of death could have been starvation since the primary food source for bull sharks resides in salt water.
In a research experiment, the bull sharks were found to be at the mouth of an estuary for the majority of the time. They stayed at the mouth of the river independent of the salinity of the water. The driving factor for a bull shark to be in fresh or salt water, however, is its age; as the bull shark ages, its tolerance for very low or high salinity increases. The majority of the newborn or very young bull sharks were found in the freshwater area, whereas the much older bull sharks were found to be in the saltwater areas, as they had developed a much better tolerance for the salinity. Reproduction is one of the reasons why adult bull sharks travel into the river—it is thought to be a physiological strategy to improve juvenile survival and a way to increase overall fitness of bull sharks. The young are not born with a high tolerance for high salinity, so they are born in fresh water and stay there until they are able to travel out.
Initially, scientists thought the sharks in Lake Nicaragua belonged to an endemic species, the Lake Nicaragua shark (Carcharhinus nicaraguensis). In 1961, following specimen comparisons, taxonomists synonymized them. Bull sharks tagged inside the lake have later been caught in the open ocean (and vice versa), with some taking as few as seven to 11 days to complete the journey.
A study of six bull sharks confined to a stagnant golf course lake in Brisbane, Australia, from 1996 to 2013 uncovered their adaptability to low-salinity environments, marking the longest recorded residency for the species under such conditions, and demonstrating their ability to live indefinitely in low-salinity aquatic environments.
Diet
The bull shark is a marine apex predator, capable of taking a variety of prey. The bull shark's diet consists mainly of bony fish and small sharks, including other bull sharks, and stingrays. Their diet can also include turtles, birds, dolphins, terrestrial mammals, crustaceans, and echinoderms. They hunt in murky waters where it is harder for the prey to see the shark coming. Bull sharks have been known to use the bump-and-bite technique to attack their prey. After the first initial contact, they continue to bite and tackle prey until the prey is unable to flee.
The bull shark is a solitary hunter, though may briefly pair with another bull shark to make hunting and tricking prey easier.
Sharks are opportunistic feeders, and the bull shark is no exception to this, as it is part of the Carcharhinus family of sharks. Normally, sharks eat in short bursts, and when food is scarce, sharks digest for a much longer period of time in order to avoid starvation. As part of their survival mechanism, bull sharks will regurgitate the food in their stomachs in order to escape from a predator. This is a distraction tactic; if the predator moves to eat the regurgitated food the bull shark can use the opportunity to escape.
Reproduction
Bull sharks mate during late summer and early autumn, often in bays and estuaries. After gestating for 12 months, a bull shark may give birth to 1 to 13 live young.
They are viviparous, born live and free-swimming. The young are about at birth. The bull shark does not rear its young; the young bull sharks are born into flat, protected areas. Coastal lagoons, river mouths, and other low-salinity estuaries are common nursery habitats.
The male bull shark is able to begin reproducing around the age of 15 years while the female cannot begin reproducing until the age of 18 years. The size of a fully matured female bull shark to produce viable eggs for fertilization seems to be 175 cm to 235 cm. The courting routine between bull sharks has not been observed in detail as of yet. The male likely bites the female on the tail until she can turn upside down and the male can copulate at that point. Mature females commonly have scratches from the mating process.
Interactions with humans
Since bull sharks often dwell in very shallow waters, are found in many types of habitats, are territorial by nature, and have no tolerance for provocation, they may be more dangerous to humans than any other species of shark. Bull sharks are one of the three shark species (along with the tiger shark and great white shark) most likely to bite humans.
One or several bull sharks may have been responsible for the Jersey Shore shark attacks of 1916. While it is a common misconception that these attacks were the inspiration for the novel Jaws, its author Peter Benchley has stated this is not the case. The speculation that bull sharks may have been responsible is based on two fatal bites occurring in brackish and fresh water.
Bull sharks have attacked swimmers around the Sydney Harbour inlets. In India, bull sharks swim up the Ganges, Bramaputra, Mahanadi, and other Indian rivers and have bitten bathers. Many of these bite incidents were attributed to the Ganges shark, Glyphis gangeticus, a critically endangered river shark species, although the sand tiger shark was also blamed during the 1960s and 1970s.
Bull sharks have also attacked humans off the coast of Florida.
Visual cues
Behavioral studies have confirmed that sharks can take visual cues to discriminate between different objects. The bull shark is able to discriminate between colors of mesh netting that is present underwater. It was found that bull sharks tended to avoid mesh netting of bright colors rather than colors that blended in with the water. Bright yellow mesh netting was found to be easily avoided when it was placed in the path of the bull shark. This was found to be the reason that sharks are attracted to bright yellow survival gear rather than ones that were painted black.
Energy conservation
In 2008, researchers tagged and recorded the movements of young bull sharks in the Caloosahatchee River estuary. They were testing to find out what determined the movement of the young bull sharks. It was found that the young bull sharks synchronously moved downriver when the environmental conditions changed. This large movement of young bull sharks were found to be moving as a response rather than other external factors such as predators. The movement was found to be directly related to the bull shark conserving energy for itself. One way the bull shark is able to conserve energy is that when the tidal flow changes, the bull shark uses the tidal flow in order to conserve energy as it moves downriver. Another way for the bull shark to conserve energy is to decrease the amount of energy needed to osmoregulate the surrounding environment.
Ecology
Humans are the biggest threat to bull sharks. Larger sharks, such as the tiger shark and great white shark, may attack them, but typically only target juveniles. Crocodiles may be a threat to bull sharks in rivers. Saltwater crocodiles have been observed preying on bull sharks in the rivers and estuaries of Northern Australia, and a Nile crocodile was reportedly sighted consuming a bull shark in South Africa.
Conservation
The bull shark is listed as a vulnerable species on the IUCN Red List and the population is currently reported as decreasing. Despite their status, the species isn't named as a protected species. Threats to the bull shark are numerous, such as getting caught in fishing nets, overfishing for their oil, skin, and meat, pollution to their habitat, and more. In many areas of the world, including Australia and South Africa, there are shark-culling measures around beaches to prevent attacks on beach-goers. Researchers tried to fix the problem of sharks getting too close to land by testing out a device called the SharkSafe Barrier™. This barrier used magnetic and visual stimuli, which utilized rows of piping to create a continuous magnetic field to deter the sharks. However, researchers concluded that the technology needs to be improved upon and tested further before it can be implemented as a reliable safety measure. Other research is being conducted to come up with conservation solutions for the bull sharks. One example is The Nature Conservancy satellite tagging sharks to track their migration and find their habitats in order to guide what areas require further protection projects.
| Biology and health sciences | Sharks | Animals |
277608 | https://en.wikipedia.org/wiki/Ploceidae | Ploceidae | Ploceidae is a family of small passerine birds, many of which are called weavers, weaverbirds, weaver finches, or bishops. These names come from the nests of intricately woven vegetation created by birds in this family. In most recent classifications, the Ploceidae are a clade that excludes some birds that have historically been placed in the family, such as some of the sparrows, but which includes the monotypic subfamily Amblyospizinae. The family is believed to have originated in the mid-Miocene. All birds of the Ploceidae are native to the Old World, most in Africa south of the Sahara, though a few live in tropical areas of Asia. A few species have been introduced outside their native range.
Taxonomy and systematics
The family Ploceidae was introduced (as Ploceïdes) by Swedish zoologist Carl Jakob Sundevall in 1836. Phylogenetic studies have shown that the family is sister to a clade containing the families Viduidae and Estrildidae Their common ancestor lived in the middle Miocene around 18 million years ago.
A 2017 molecular phylogenetic study by Thilina de Silva and collaborators, as well as an expanded study by the same group published in 2019 have indicated that the genus Ploceus as currently defined is polyphyletic. A cladogram based on these results is shown below.
Genera
The family includes 15 genera with a total of 122 species. For more detail, see list of Ploceidae species.
Description
The males of many species in this family are brightly coloured, usually in red or yellow and black. Some species show variation in colour only in the breeding season. These are seed-eating birds with rounded conical bills.
Distribution and habitat
The weaverbird colonies may be found close to bodies of water.
Behaviour and ecology
Weavers are named for their elaborately woven nests. The nests vary in size, shape, material used, and construction techniques from species to species. Materials used for building nests include fine leaf fibers, grass, and twigs. Many species weave very fine nests using thin strands of leaf fiber, though some, like the buffalo-weavers, form massive untidy stick nests in their colonies, which may have spherical woven nests within. The sociable weavers of Africa build apartment-house nests, in which 100 to 300 pairs have separate flask-shaped chambers entered by tubes at the bottom. The sparrow weavers live in family units that employ cooperative breeding.
Most species weave nests that have narrow entrances, facing downward.
Many weaver species are gregarious and breed colonially. The birds build their nests together for protection, often several to a branch. Usually the male birds weave the nests and use them as a form of display to lure prospective females.
Relationship to humans
They sometimes cause crop damage, notably the red-billed quelea, reputed to be the world's most numerous bird.
Gallery
| Biology and health sciences | Passerida | null |
277641 | https://en.wikipedia.org/wiki/Strength%20of%20materials | Strength of materials | The field of strength of materials (also called mechanics of materials) typically refers to various methods of calculating the stresses and strains in structural members, such as beams, columns, and shafts. The methods employed to predict the response of a structure under loading and its susceptibility to various failure modes takes into account the properties of the materials such as its yield strength, ultimate strength, Young's modulus, and Poisson's ratio. In addition, the mechanical element's macroscopic properties (geometric properties) such as its length, width, thickness, boundary constraints and abrupt changes in geometry such as holes are considered.
The theory began with the consideration of the behavior of one and two dimensional members of structures, whose states of stress can be approximated as two dimensional, and was then generalized to three dimensions to develop a more complete theory of the elastic and plastic behavior of materials. An important founding pioneer in mechanics of materials was Stephen Timoshenko.
Definition
In the mechanics of materials, the strength of a material is its ability to withstand an applied load without failure or plastic deformation. The field of strength of materials deals with forces and deformations that result from their acting on a material. A load applied to a mechanical member will induce internal forces within the member called stresses when those forces are expressed on a unit basis. The stresses acting on the material cause deformation of the material in various manners including breaking them completely. Deformation of the material is called strain when those deformations too are placed on a unit basis.
The stresses and strains that develop within a mechanical member must be calculated in order to assess the load capacity of that member. This requires a complete description of the geometry of the member, its constraints, the loads applied to the member and the properties of the material of which the member is composed. The applied loads may be axial (tensile or compressive), or rotational (strength shear). With a complete description of the loading and the geometry of the member, the state of stress and state of strain at any point within the member can be calculated. Once the state of stress and strain within the member is known, the strength (load carrying capacity) of that member, its deformations (stiffness qualities), and its stability (ability to maintain its original configuration) can be calculated.
The calculated stresses may then be compared to some measure of the strength of the member such as its material yield or ultimate strength. The calculated deflection of the member may be compared to deflection criteria that are based on the member's use. The calculated buckling load of the member may be compared to the applied load. The calculated stiffness and mass distribution of the member may be used to calculate the member's dynamic response and then compared to the acoustic environment in which it will be used.
Material strength refers to the point on the engineering stress–strain curve (yield stress) beyond which the material experiences deformations that will not be completely reversed upon removal of the loading and as a result, the member will have a permanent deflection. The ultimate strength of the material refers to the maximum value of stress reached. The fracture strength is the stress value at fracture (the last stress value recorded).
Types of loadings
Transverse loadings – Forces applied perpendicular to the longitudinal axis of a member. Transverse loading causes the member to bend and deflect from its original position, with internal tensile and compressive strains accompanying the change in curvature of the member. Transverse loading also induces shear forces that cause shear deformation of the material and increase the transverse deflection of the member.
Axial loading – The applied forces are collinear with the longitudinal axis of the member. The forces cause the member to either stretch or shorten.
Torsional loading – Twisting action caused by a pair of externally applied equal and oppositely directed force couples acting on parallel planes or by a single external couple applied to a member that has one end fixed against rotation.
Stress terms
Uniaxial stress is expressed by
where F is the force acting on an area A. The area can be the undeformed area or the deformed area, depending on whether engineering stress or true stress is of interest.
Compressive stress (or compression) is the stress state caused by an applied load that acts to reduce the length of the material (compression member) along the axis of the applied load; it is, in other words, a stress state that causes a squeezing of the material. A simple case of compression is the uniaxial compression induced by the action of opposite, pushing forces. Compressive strength for materials is generally higher than their tensile strength. However, structures loaded in compression are subject to additional failure modes, such as buckling, that are dependent on the member's geometry.
Tensile stress is the stress state caused by an applied load that tends to elongate the material along the axis of the applied load, in other words, the stress caused by pulling the material. The strength of structures of equal cross-sectional area loaded in tension is independent of shape of the cross-section. Materials loaded in tension are susceptible to stress concentrations such as material defects or abrupt changes in geometry. However, materials exhibiting ductile behaviour (many metals for example) can tolerate some defects while brittle materials (such as ceramics and some steels) can fail well below their ultimate material strength.
Shear stress is the stress state caused by the combined energy of a pair of opposing forces acting along parallel lines of action through the material, in other words, the stress caused by faces of the material sliding relative to one another. An example is cutting paper with scissors or stresses due to torsional loading.
Stress parameters for resistance
Material resistance can be expressed in several mechanical stress parameters. The term material strength is used when referring to mechanical stress parameters. These are physical quantities with dimension homogeneous to pressure and force per unit surface. The traditional measure unit for strength are therefore MPa in the International System of Units, and the psi between the United States customary units. Strength parameters include: yield strength, tensile strength, fatigue strength, crack resistance, and other parameters.
Yield strength is the lowest stress that produces a permanent deformation in a material. In some materials, like aluminium alloys, the point of yielding is difficult to identify, thus it is usually defined as the stress required to cause 0.2% plastic strain. This is called a 0.2% proof stress.
Compressive strength is a limit state of compressive stress that leads to failure in a material in the manner of ductile failure (infinite theoretical yield) or brittle failure (rupture as the result of crack propagation, or sliding along a weak plane – see shear strength).
Tensile strength or ultimate tensile strength is a limit state of tensile stress that leads to tensile failure in the manner of ductile failure (yield as the first stage of that failure, some hardening in the second stage and breakage after a possible "neck" formation) or brittle failure (sudden breaking in two or more pieces at a low-stress state). The tensile strength can be quoted as either true stress or engineering stress, but engineering stress is the most commonly used.
Fatigue strength is a more complex measure of the strength of a material that considers several loading episodes in the service period of an object, and is usually more difficult to assess than the static strength measures. Fatigue strength is quoted here as a simple range (). In the case of cyclic loading it can be appropriately expressed as an amplitude usually at zero mean stress, along with the number of cycles to failure under that condition of stress.
Impact strength is the capability of the material to withstand a suddenly applied load and is expressed in terms of energy. Often measured with the Izod impact strength test or Charpy impact test, both of which measure the impact energy required to fracture a sample. Volume, modulus of elasticity, distribution of forces, and yield strength affect the impact strength of a material. In order for a material or object to have a high impact strength, the stresses must be distributed evenly throughout the object. It also must have a large volume with a low modulus of elasticity and a high material yield strength.
Strain parameters for resistance
Deformation of the material is the change in geometry created when stress is applied (as a result of applied forces, gravitational fields, accelerations, thermal expansion, etc.). Deformation is expressed by the displacement field of the material.
Strain, or reduced deformation, is a mathematical term that expresses the trend of the deformation change among the material field. Strain is the deformation per unit length. In the case of uniaxial loading the displacement of a specimen (for example, a bar element) lead to a calculation of strain expressed as the quotient of the displacement and the original length of the specimen. For 3D displacement fields it is expressed as derivatives of displacement functions in terms of a second-order tensor (with 6 independent elements).
Deflection is a term to describe the magnitude to which a structural element is displaced when subject to an applied load.
Stress–strain relations
Elasticity is the ability of a material to return to its previous shape after stress is released. In many materials, the relation between applied stress is directly proportional to the resulting strain (up to a certain limit), and a graph representing those two quantities is a straight line.
The slope of this line is known as Young's modulus, or the "modulus of elasticity". The modulus of elasticity can be used to determine the stress–strain relationship in the linear-elastic portion of the stress–strain curve. The linear-elastic region is either below the yield point, or if a yield point is not easily identified on the stress–strain plot it is defined to be between 0 and 0.2% strain, and is defined as the region of strain in which no yielding (permanent deformation) occurs.
Plasticity or plastic deformation is the opposite of elastic deformation and is defined as unrecoverable strain. Plastic deformation is retained after the release of the applied stress. Most materials in the linear-elastic category are usually capable of plastic deformation. Brittle materials, like ceramics, do not experience any plastic deformation and will fracture under relatively low strain, while ductile materials such as metallics, lead, or polymers will plastically deform much more before a fracture initiation.
Consider the difference between a carrot and chewed bubble gum. The carrot will stretch very little before breaking. The chewed bubble gum, on the other hand, will plastically deform enormously before finally breaking.
Design terms
Ultimate strength is an attribute related to a material, rather than just a specific specimen made of the material, and as such it is quoted as the force per unit of cross section area (N/m2). The ultimate strength is the maximum stress that a material can withstand before it breaks or weakens. For example, the ultimate tensile strength (UTS) of AISI 1018 Steel is 440 MPa. In Imperial units, the unit of stress is given as lbf/in2 or pounds-force per square inch. This unit is often abbreviated as psi. One thousand psi is abbreviated ksi.
A factor of safety is a design criteria that an engineered component or structure must achieve. , where FS: the factor of safety, Rf The applied stress, and F: ultimate allowable stress (psi or MPa)
Margin of Safety is the common method for design criteria. It is defined MS = Pu/P − 1.
For example, to achieve a factor of safety of 4, the allowable stress in an AISI 1018 steel component can be calculated to be = 440/4 = 110 MPa, or = 110×106 N/m2. Such allowable stresses are also known as "design stresses" or "working stresses".
Design stresses that have been determined from the ultimate or yield point values of the materials give safe and reliable results only for the case of static loading. Many machine parts fail when subjected to a non-steady and continuously varying loads even though the developed stresses are below the yield point. Such failures are called fatigue failure. The failure is by a fracture that appears to be brittle with little or no visible evidence of yielding. However, when the stress is kept below "fatigue stress" or "endurance limit stress", the part will endure indefinitely. A purely reversing or cyclic stress is one that alternates between equal positive and negative peak stresses during each cycle of operation. In a purely cyclic stress, the average stress is zero. When a part is subjected to a cyclic stress, also known as stress range (Sr), it has been observed that the failure of the part occurs after a number of stress reversals (N) even if the magnitude of the stress range is below the material's yield strength. Generally, higher the range stress, the fewer the number of reversals needed for failure.
Failure theories
There are four failure theories: maximum shear stress theory, maximum normal stress theory, maximum strain energy theory, and maximum distortion energy theory (von Mises criterion of failure). Out of these four theories of failure, the maximum normal stress theory is only applicable for brittle materials, and the remaining three theories are applicable for ductile materials.
Of the latter three, the distortion energy theory provides the most accurate results in a majority of the stress conditions. The strain energy theory needs the value of Poisson's ratio of the part material, which is often not readily available. The maximum shear stress theory is conservative. For simple unidirectional normal stresses all theories are equivalent, which means all theories will give the same result.
Maximum shear stress theory postulates that failure will occur if the magnitude of the maximum shear stress in the part exceeds the shear strength of the material determined from uniaxial testing.
Maximum normal stress theory postulates that failure will occur if the maximum normal stress in the part exceeds the ultimate tensile stress of the material as determined from uniaxial testing. This theory deals with brittle materials only. The maximum tensile stress should be less than or equal to ultimate tensile stress divided by factor of safety. The magnitude of the maximum compressive stress should be less than ultimate compressive stress divided by factor of safety.
Maximum strain energy theory postulates that failure will occur when the strain energy per unit volume due to the applied stresses in a part equals the strain energy per unit volume at the yield point in uniaxial testing.
Maximum distortion energy theory, also known as maximum distortion energy theory of failure or von Mises–Hencky theory. This theory postulates that failure will occur when the distortion energy per unit volume due to the applied stresses in a part equals the distortion energy per unit volume at the yield point in uniaxial testing. The total elastic energy due to strain can be divided into two parts: one part causes change in volume, and the other part causes a change in shape. Distortion energy is the amount of energy that is needed to change the shape.
Fracture mechanics was established by Alan Arnold Griffith and George Rankine Irwin. This important theory is also known as numeric conversion of toughness of material in the case of crack existence.
A material's strength depends on its microstructure. The engineering processes to which a material is subjected can alter its microstructure. Strengthening mechanisms that alter the strength of a material include work hardening, solid solution strengthening, precipitation hardening, and grain boundary strengthening.
Strengthening mechanisms are accompanied by the caveat that some other mechanical properties of the material may degenerate in an attempt to make a material stronger. For example, in grain boundary strengthening, although yield strength is maximized with decreasing grain size, ultimately, very small grain sizes make the material brittle. In general, the yield strength of a material is an adequate indicator of the material's mechanical strength. Considered in tandem with the fact that the yield strength is the parameter that predicts plastic deformation in the material, one can make informed decisions on how to increase the strength of a material depending on its microstructural properties and the desired end effect. Strength is expressed in terms of the limiting values of the compressive stress, tensile stress, and shear stresses that would cause failure. The effects of dynamic loading are probably the most important practical consideration of the theory of elasticity, especially the problem of fatigue. Repeated loading often initiates cracks, which grow until failure occurs at the corresponding residual strength of the structure. Cracks always start at a stress concentrations especially changes in cross-section of the product or defects in manufacturing, near holes and corners at nominal stress levels far lower than those quoted for the strength of the material.
| Physical sciences | Solid mechanics | null |
277663 | https://en.wikipedia.org/wiki/Open-source%20software | Open-source software | Open-source software (OSS) is computer software that is released under a license in which the copyright holder grants users the rights to use, study, change, and distribute the software and its source code to anyone and for any purpose. Open-source software may be developed in a collaborative, public manner. Open-source software is a prominent example of open collaboration, meaning any capable user is able to participate online in development, making the number of possible contributors indefinite. The ability to examine the code facilitates public trust in the software.
Open-source software development can bring in diverse perspectives beyond those of a single company. A 2024 estimate of the value of open-source software to firms is $8.8 trillion, as firms would need to spend 3.5 times the amount they currently do without the use of open source software.
Open-source code can be used for studying and allows capable end users to adapt software to their personal needs in a similar way user scripts and custom style sheets allow for web sites, and eventually publish the modification as a fork for users with similar preferences, and directly submit possible improvements as pull requests.
Definitions
The Open Source Initiative's (OSI) definition is recognized by several governments internationally as the standard or de facto definition. OSI uses The Open Source Definition to determine whether it considers a software license open source. The definition was based on the Debian Free Software Guidelines, written and adapted primarily by Perens. Perens did not base his writing on the "four freedoms" from the Free Software Foundation (FSF), which were only widely available later.
Under Perens' definition, open source is a broad software license that makes source code available to the general public with relaxed or non-existent restrictions on the use and modification of the code. It is an explicit "feature" of open source that it puts very few restrictions on the use or distribution by any organization or user, in order to enable the rapid evolution of the software.
According to Feller et al. (2005), the terms "free software" and "open-source software" should be applied to any "software products distributed under terms that allow users" to use, modify, and redistribute the software "in any manner they see fit, without requiring that they pay the author(s) of the software a royalty or fee for engaging in the listed activities."
Despite initially accepting it, Richard Stallman of the FSF now flatly opposes the term "Open Source" being applied to what they refer to as "free software". Although he agrees that the two terms describe "almost the same category of software", Stallman considers equating the terms incorrect and misleading. Stallman also opposes the professed pragmatism of the Open Source Initiative, as he fears that the free software ideals of freedom and community are threatened by compromising on the FSF's idealistic standards for software freedom. The FSF considers free software to be a subset of open-source software, and Richard Stallman explained that DRM software, for example, can be developed as open source, despite that it does not give its users freedom (it restricts them), and thus does not qualify as free software.
Open-source software development
Development model
In his 1997 essay The Cathedral and the Bazaar, open-source influential contributor Eric S. Raymond suggests a model for developing OSS known as the bazaar model. Raymond likens the development of software by traditional methodologies to building a cathedral, with careful isolated work by individuals or small groups. He suggests that all software should be developed using the bazaar style, with differing agendas and approaches.
In the traditional model of development, which he called the cathedral model, development takes place in a centralized way. Roles are clearly defined. Roles include people dedicated to designing (the architects), people responsible for managing the project, and people responsible for implementation. Traditional software engineering follows the cathedral model.
The bazaar model, however, is different. In this model, roles are not clearly defined. Some proposed characteristics of software developed using the bazaar model should exhibit the following patterns:
Users should be treated as co-developers: The users are treated like co-developers and so they should have access to the source code of the software. Furthermore, users are encouraged to submit additions to the software, code fixes for the software, bug reports, documentation, etc. Having more co-developers increases the rate at which the software evolves. Linus's law states that given enough eyeballs all bugs are shallow. This means that if many users view the source code, they will eventually find all bugs and suggest how to fix them. Some users have advanced programming skills, and furthermore, each user's machine provides an additional testing environment. This new testing environment offers the ability to find and fix a new bug.
Early releases: The first version of the software should be released as early as possible so as to increase one's chances of finding co-developers early.
Frequent integration: Code changes should be integrated (merged into a shared code base) as often as possible so as to avoid the overhead of fixing a large number of bugs at the end of the project life cycle. Some open-source projects have nightly builds where integration is done automatically.
Several versions: There should be at least two versions of the software. There should be a buggier version with more features and a more stable version with fewer features. The buggy version (also called the development version) is for users who want the immediate use of the latest features and are willing to accept the risk of using code that is not yet thoroughly tested. The users can then act as co-developers, reporting bugs and providing bug fixes.
High modularization: The general structure of the software should be modular allowing for parallel development on independent components.
Dynamic decision-making structure: There is a need for a decision-making structure, whether formal or informal, that makes strategic decisions depending on changing user requirements and other factors. Compare with extreme programming.
The process of Open source development begins with a requirements elicitation where developers consider if they should add new features or if a bug needs to be fixed in their project. This is established by communicating with the OSS community through avenues such as bug reporting and tracking or mailing lists and project pages. Next, OSS developers select or are assigned to a task and identify a solution. Because there are often many different possible routes for solutions in OSS, the best solution must be chosen with careful consideration and sometimes even peer feedback. The developer then begins to develop and commit the code. The code is then tested and reviewed by peers. Developers can edit and evolve their code through feedback from continuous integration. Once the leadership and community are satisfied with the whole project, it can be partially released and user instruction can be documented. If the project is ready to be released, it is frozen, with only serious bug fixes or security repairs occurring. Finally, the project is fully released and only changed through minor bug fixes.
Advantages
Open source implementation of a standard can increase adoption of that standard. This creates developer loyalty as developers feel empowered and have a sense of ownership of the end product.
Moreover, lower costs of marketing and logistical services are needed for OSS. OSS can be a tool to promote a company's image, including its commercial products. The OSS development approach has helped produce reliable, high quality software quickly and inexpensively.
Open source development offers the potential to quicken innovation and create of social value. In France for instance, a policy that incentivized government to favor free open-source software increased to nearly 600,000 OSS contributions per year, generating social value by increasing the quantity and quality of open-source software. This policy also led to an estimated increase of up to 18% of tech startups and a 14% increase in the number of people employed in the IT sector.
OSS can be highly reliable when it has thousands of independent programmers testing and fixing bugs of the software. Open source is not dependent on the company or author that originally created it. Even if the company fails, the code continues to exist and be developed by its users.
OSS is flexible because modular systems allow programmers to build custom interfaces, or add new abilities to it and it is innovative since open-source programs are the product of collaboration among a large number of different programmers. The mix of divergent perspectives, corporate objectives, and personal goals speeds up innovation.
Moreover, free software can be developed in accordance with purely technical requirements. It does not require thinking about commercial pressure that often degrades the quality of the software. Commercial pressures make traditional software developers pay more attention to customers' requirements than to security requirements, since such features are somewhat invisible to the customer.
Development tools
In open-source software development, tools are used to support the development of the product and the development process itself.
Version control systems such as Centralized Version control system (CVCS) and the distributed version control system (DVCS) are examples of tools, often open source, that help manage the source code files and the changes to those files for a software project in order to foster collaboration. CVCS are centralized with a central repository while DVCS are decentralized and have a local repository for every user. concurrent versions system (CVS) and later Subversion (SVN) and Git are examples of CVCS. The repositories are hosted and published on source-code-hosting facilities such as GitHub.
Open-source projects use utilities such as issue trackers to organize open-source software development. Commonly used bug trackers include Bugzilla and Redmine.
Tools such as mailing lists and IRC provide means of coordination and discussion of bugs among developers. Project web pages, wiki pages, roadmap lists and newsgroups allow for the distribution of project information that focuses on end users.
Opportunities for participation
Contributing
The basic roles OSS participants can fall into multiple categories, beginning with leadership at the center of the project who have control over its execution. Next are the core contributors with a great deal of experience and authority in the project who may guide the other contributors. Non-core contributors have less experience and authority, but regularly contribute and are vital to the project's development. New contributors are the least experienced but with mentorship and guidance can become regular contributors.
Some possible ways of contributing to open-source software include such roles as programming, user interface design and testing, web design, bug triage, accessibility design and testing, UX design, code testing, and security review and testing. However, there are several ways of contributing to OSS projects even without coding skills. For example, some less technical ways of participating are documentation writing and editing, translation, project management, event organization and coordination, marketing, release management, community management, and public relations and outreach.
Funding is absolutely another terrific way that individuals and organizations choose to contribute to open source projects. Groups like Open Collective provide a means for individuals to contribute monthly to supporting their favorite projects. Organizations like the Sovereign Tech Fund is able to contribute to millions to supporting the tools the German Government uses. The National Science Foundation established a Pathways to Enable Open-Source Ecosystems (POSE) program to support open source innovation.
Industry participation
The adoption of open-source software by industry is increasing over time. OSS is popular in several industries such as telecommunications, aerospace, healthcare, and media & entertainment due to the benefits it provides. Adoption of OSS is more likely in larger organizations and is dependent on the company's IT usage, operating efficiencies, and the productivity of employees.
Industries are likely to use OSS due to back-office functionality, sales support, research and development, software features, quick deployment, portability across platforms and avoidance of commercial license management. Additionally, lower cost for hardware and ownership are also important benefits.
Prominent organizations
Organizations that contribute to the development and expansions of free and open-source software movements exist all over the world. These organizations are dedicated to goals such as teaching and spreading technology. As listed by a former vice president of the Open Source Initiative, some American organizations include the Free Software Foundation, Software Freedom Conservancy, the Open Source Initiative and Software in the Public Interest. Within Europe some notable organizations are Free Software Foundation Europe, open-source projects EU (OSP) and OpenForum Europe (OFE). One Australian organization is Linux Australia while Asia has Open source Asia and FOSSAsia. Free and open source software for Africa (FOSSFA) and OpenAfrica are African organizations and Central and South Asia has such organizations as FLISOL and GRUP de usuarios de software libre Peru. Outside of these, many more organizations dedicated to the advancement of open-source software exist.
Legal and economic issues
Licensing
FOSS products are generally licensed under two types of licenses: permissive licensing and copyleft licensing. Both of these types of licenses are different than proprietary licensing in that they can allow more users access to the software and allow for the creation of derivative works as specified by the terms of the specific license, as each license has its own rules. Permissive licenses allow recipients of the software to implement the author's copyright rights without having to use the same license for distribution. Examples of this type of license include the BSD, MIT, and Apache licenses. Copyleft licenses are different in that they require recipients to use the same license for at least some parts of the distribution of their works. Strong copyleft licenses require all derivative works to use the same license while weak copyleft licenses require the use of the same license only under certain conditions. Examples of this type of license include the GNU family of licenses, and the MPL and EPL licenses. The similarities between these two categories of licensing include that they provide a broad grant of copyright rights, require that recipients preserve copyright notices, and that a copy of the license is provided to recipients with the code.
One important legal precedent for open-source software was created in 2008, when the Jacobson v Katzer case enforced terms of the Artistic license, including attribution and identification of modifications. The ruling of this case cemented enforcement under copyright law when the conditions of the license were not followed. Because of the similarity of the Artistic license to other open-source software licenses, the ruling created a precedent that applied widely.
Examples of free-software license / open-source licenses include Apache licenses, BSD licenses, GNU General Public Licenses, GNU Lesser General Public License, MIT License, Eclipse Public License and Mozilla Public License.
Legal issues
Several gray areas exist within software regulation that have great impact on open-source software, such as if software is a good or service, what can be considered a modification, governance through contract vs license, ownership and right of use. While there have been developments on these issues, they often lead to even more questions. The existence of these uncertainties in regulation has a negative impact on industries involved in technologies as a whole.
Within the legal history of software as a whole, there was much debate on whether to protect it as intellectual property under patent law, copyright law or establishing a unique regulation. Ultimately, copyright law became the standard with computer programs being considered a form of literary work, with some tweaks of unique regulation.
Software is generally considered source code and object code, with both being protectable, though there is legal variety in this definition. Some jurisdictions attempt to expand or reduce this conceptualization for their own purposes. For example, The European Court of Justice defines a computer program as not including the functionality of a program, the programing language, or the format of data files. By limiting protections of the different aspects of software, the law favors an open-source approach to software use. The US especially has an open approach to software, with most open-source licenses originating there. However, this has increased the focus on patent rights within these licenses, which has seen backlash from the OSS community, who prefer other forms of IP protection.
Another issue includes technological protection measures (TPM) and digital rights management (DRM) techniques which were internationally legally recognized and protected in the 1996 World Intellectual Property Organization (WIPO) Treaty. Open source software proponents disliked these technologies as they constrained end-users potentially beyond copyright law. Europe responded to such complaints by putting TPM under legal controls, representing a victory for OSS supporters.
Economic/business implications
In open-source communities, instead of owning the software produced, the producer owns the development of the evolving software. In this way, the future of the software is open, making ownership or intellectual property difficult within OSS. Licensing and branding can prevent others from stealing it, preserving its status as a public good. Open source software can be considered a public good as it is available to everyone and does not decrease in value for others when downloaded by one person. Open source software is unique in that it becomes more valuable as it is used and contributed to, instead of diminishing the resource. This is explained by concepts such as investment in reputation and network effects.
The economic model of open-source software can be explained as developers contribute work to projects, creating public benefits. Developers choose projects based on the perceived benefits or costs, such as improved reputation or value of the project. The motivations of developers can come from many different places and reasons, but the important takeaway is that money is not the only or even most important incentivization.
Because economic theory mainly focuses on the consumption of scarce resources, the OSS dynamic can be hard to understand. In OSS, producers become consumers by reaping the rewards of contributing to a project. For example, a developer becomes well regarded by their peers for a successful contribution to an OSS project. The social benefits and interactions of OSS are difficult to account for in economic models as well. Furthermore, the innovation of technology creates constantly changing value discussions and outlooks, making economic model unable to predict social behavior.
Although OSS is theoretically challenging in economic models, it is explainable as a sustainable social activity that requires resources. These resources include time, money, technology and contributions. Many developers have used technology funded by organizations such as universities and governments, though these same organizations benefit from the work done by OSS. As OSS grows, hybrid systems containing OSS and proprietary systems are becoming more common.
Throughout the mid 2000s, more and more tech companies have begun to use OSS. For example, Dell's move of selling computers with GNU/Linux already installed. Microsoft itself has launched a Linux-based operating system despite previous animosity with the OSS movement. Despite these developments, these companies tend to only use OSS for certain purposes, leading to worries that OSS is being taken advantage of by corporations and not given anything in return.
Government uses
While many governments are interested in implementing and promoting open-source software due to the many benefits provided, a huge issue to be considered is cybersecurity. While accidental vulnerabilities are possible, so are attacks by outside agents. Because of these fears, governmental interest in contributing to the governance of software has become more prominent. However, these are the broad strokes of the issue, with each country having their own specific politicized interactions with open-source software and their goals for its implementation. For example, the United States has focused on national security in regard to open-source software implementation due to the perceived threat of the increase of open-source software activity in countries like China and Russia, with the Department of Defense considering multiple criteria for using OSS. These criteria include: if it comes from and is maintained by trusted sources, whether it will continue to be maintained, if there are dependencies on sub-components in the software, component security and integrity, and foreign governmental influence.
Another issue for governments in regard to open source is their investments in technologies such as operating systems, semiconductors, cloud, and artificial intelligence. These technologies all have implications for global cooperation, again opening up security issues and political consequences. Many countries have to balance technological innovation with technological dependence in these partnerships. For example, after China's open-source dependent company Huawei was prevented from using Google's Android system in 2019, they began to create their own alternative operating system: Harmony OS.
Germany recently established a Sovereign Tech Fund, to help support the governance and maintenance of the software that they use.
Open software movement
History
In the early days of computing, such as the 1950s and into the 1960s, programmers and developers shared software to learn from each other and evolve the field of computing. For example, Unix included the operating system source code for users. Eventually, the commercialization of software in the years 1970–1980 began to prevent this practice. However, academics still often developed software collaboratively.
In response, the open-source movement was born out of the work of skilled programmer enthusiasts, widely referred to as hackers or hacker culture. One of these enthusiasts, Richard Stallman, was a driving force behind the free software movement, which would later allow for the open-source movement. In 1984, he resigned from MIT to create a free operating system, GNU, after the programmer culture in his lab was stifled by proprietary software preventing source code from being shared and improved upon. GNU was UNIX compatible, meaning that the programmer enthusiasts would still be familiar with how it worked. However, it quickly became apparent that there was some confusion with the label Stallman had chosen of free software, which he described as free as in free speech, not free beer, referring to the meaning of free as freedom rather than price. He later expanded this concept of freedom to the four essential freedoms. Through GNU, open-source norms of incorporating others' source code, community bug fixes and suggestions of code for new features appeared. In 1985, Stallman founded the Free Software Foundation (FSF) to promote changes in software and to help write GNU. In order to prevent his work from being used in proprietary software, Stallman created the concept of copyleft, which allowed the use of his work by anyone, but under specific terms. To do this, he created the GNU General Public License (GNU GPL) in 1989, which was updated in 1991. In 1991, GNU was combined with the Linux kernel written by Linus Torvalds, as a kernel was missing in GNU. The operating system is now usually referred to as Linux. Throughout this whole period, there were many other free software projects and licenses around at the time, all with different ideas of what the concept of free software was and should be, as well as the morality of proprietary software, such as Berkeley Software Distribution, TeX, and the X Window System.
As free software developed, the Free Software Foundation began to look how to bring free software ideas and perceived benefits to the commercial software industry. It was concluded that FSF's social activism was not appealing to companies and they needed a way to rebrand the free software movement to emphasize the business potential of sharing and collaborating on software source code. The term open source was suggested by Christine Peterson in 1998 at a meeting of supporters of free software. Many in the group felt the name free software was confusing to newcomers and holding back industry interest and they readily accepted the new designation of open source, creating the Open Source Initiative (OSI) and the OSI definition of what open source software is. The Open Source Initiative's (OSI) definition is now recognized by several governments internationally as the standard or de facto definition. The definition was based on the Debian Free Software Guidelines, written and adapted primarily by Bruce Perens. The OSI definition differed from the free software definition in that it allows the inclusion of proprietary software and allows more liberties in its licensing. Some, such as Stallman, agree more with the original concept of free software as a result because it takes a strong moral stance against proprietary software, through there is much overlap between the two movements in terms of the operation of the software.
While the Open Source Initiative sought to encourage the use of the new term and evangelize the principles it adhered to, commercial software vendors found themselves increasingly threatened by the concept of freely distributed software and universal access to an application's source code, with an executive of Microsoft calling open source an intellectual property destroyer in 2001. However, while free and open-source software (FOSS) has historically played a role outside of mainstream private software development, companies as large as Microsoft have begun to develop official open source presences on the Internet. IBM, Oracle, and State Farm are just a few of the companies with a serious public stake in today's competitive open source market, marking a significant shift in the corporate philosophy concerning the development of FOSS.
Future
The future of the open source software community, and the free software community by extension, has become successful if not confused about what it stands for. For example, Android and Ubuntu are examples milestones of success in the open source software rise to prominence from the sidelines of technological innovation as it existed in the early 2000s. However, some in the community consider them failures in their representation of OSS due to issues such as the downplaying of the OSS center of Android by Google and its partners, the use of an Apache license that allowed forking and resulted in a loss of opportunities for collaboration within Android, the prioritization of convenience over freedom in Ubuntu, and features within Ubuntu that track users for marketing purposes.
The use of OSS has become more common in business with 78% of companies reporting that they run all or part of their operations on FOSS. The popularity of OSS has risen to the point that Microsoft, a once detractor of OSS, has included its use in their systems. However, this success has raised concerns that will determine the future of OSS as the community must answer questions such as what OSS is, what should it be, and what should be done to protect it, if it even needs protecting. All in all, while the free and open source revolution has slowed to a perceived equilibrium in the market place, that does not mean it is over as many theoretical discussions must take place to determine its future.
Comparisons with other software licensing/development models
Closed source / proprietary software
Open source software differs from proprietary software in that it is publicly available, the license requires no fees, modifications and distributions are allowed under license specifications. All of this works to prevent a monopoly on any OSS product, which is a goal of proprietary software. Proprietary software limits their customers' choices to either committing to using that software, upgrading it or switching to other software, forcing customers to have their software preferences impacted by their monetary cost. The ideal case scenario for the proprietary software vendor would be a lock-in, where the customer does not or cannot switch software due to these costs and continues to buy products from that vendor.
Within proprietary software, bug fixes can only be provided by the vendor, moving platforms requires another purchase and the existence of the product relies on the vendor, who can discontinue it at any point. Additionally, proprietary software does not provide its source code and cannot be altered by users. For businesses, this can pose a security risk and source of frustration, as they cannot specialize the product to their needs, and there may be hidden threats or information leaks within the software that they cannot access or change.
Free software
Under OSI's definition, open source is a broad software license that makes source code available to the general public with relaxed or non-existent restrictions on the use and modification of the code. It is an explicit feature of open source that it puts very few restrictions on the use or distribution by any organization or user, in order to enable the rapid evolution of the software.
Richard Stallman, leader of the Free software movement and member of the free software foundation opposes the term open source being applied to what they refer to as free software. Although he agrees that the two terms describe almost the same category of software, Stallman considers equating the terms incorrect and misleading. He believes that the main difference is that by choosing one term over the other lets others know about what one's goals are: development (open source) or a social stance (free software). Nevertheless, there is significant overlap between open source software and free software. Stallman also opposes the professed pragmatism of the Open Source Initiative, as he fears that the free software ideals of freedom and community are threatened by compromising on the FSF's idealistic standards for software freedom. The FSF considers free software to be a subset of open-source software, and Richard Stallman explained that DRM software, for example, can be developed as open source, despite how it restricts its users, and thus does not qualify as free software.
The FSF said that the term open source fosters an ambiguity of a different kind such that it confuses the mere availability of the source with the freedom to use, modify, and redistribute it. On the other hand, the term free software was criticized for the ambiguity of the word free, which was seen as discouraging for business adoption, and for the historical ambiguous usage of the term.
Developers have used the alternative terms Free and Open Source Software (FOSS), or Free/Libre and Open Source Software (FLOSS), consequently, to describe open-source software that is also free software.
Source-available software
Software can be distributed with source code, which is a code that is readable. Software is source available when this source code is available to be seen. However to be source available or FOSS, the source code does not need to be accessible to all, just the users of that software. While all FOSS software is source available because this is a requirement made by the Open Source Definition, not all source available software is FOSS. For example, if the software does not meet other aspects of the Open Source Definition such as permitted modification or redistribution, even if the source code is available, the software is not FOSS.
Open-sourcing
A recent trend within software companies is open sourcing, or transitioning their previous proprietary software into open source software through releasing it under an open-source license. Examples of companies who have done this are Google, Microsoft and Apple. Additionally, open sourcing can refer to programming open source software or installing open source software. Open sourcing can be beneficial in multiple ways, such as attracting more external contributors who bring new perspectives and problem solving capabilities. The downsides of open sourcing include the work that has to be done to maintaining the new community, such as making the base code easily understandable, setting up communication channels for new developers and creating documentation to allow new developers to easily join. However, a review of several open sourced projects found that although a newly open sourced project attracts many newcomers, a great amount are likely to soon leave the project and their forks are also likely to not be impactful.
Other
Other concepts that may share some similarities to open source are shareware, public domain software, freeware, and software viewers/readers that are freely available but do not provide source code. However, these differ from open source software in access to source code, licensing, copyright and fees.
Society and culture
Demographics
Despite being able to collaborate internationally, open source software contributors were found to mostly be located in large clusters such as Silicon Valley that largely collaborate within themselves. Possible reasons for this phenomenon may be that the OSS contributor demographic largely works in software, meaning that the OSS geographic location is closely related to that dispersion and collaborations could be encouraged through work and social networks. Code acceptance can be impacted by status within these social network clusters, creating unfair predispositions in code acceptance based on location. Barriers to international collaboration also include linguistic or cultural differences. Furthermore, each country has been shown to have a higher acceptance rate for code from contributors within their country except India, indicating a bias for culturally similar collaborators.
In 2021, the countries with the highest open source software contributions included the United States, China, Germany, India, and the UK, in that order. The counties with the highest OSS developers per capita from a study in 2021 include, in order, Iceland, Switzerland, Norway, Sweden, and Finland, while in 2008 the countries with top amount of estimated contributors in SourceForge were the United States, Germany, United Kingdom, Canada and France. Though there have been several studies done on the distribution and contributions of OSS developers, this is still an open field that can be measured in several different ways. For instance, Information and communication technology participation, population, wealth and proportion of access to the internet have been shown to be correlated with OSS contributions.
Although gender diversity has been found to enhance team productivity, women still face biases while contributing to open source software projects when their gender is identifiable. In 2002, only 1.5% of international open-source software developers were women, while women made up 28% of tech industry roles, demonstrating their low representation in the software field. Despite OSS contributions having no prerequisites, this gender bias may continue to exist due to the common belief of contributors that gender should not matter, and the quality of code should be the only consideration for code acceptance, preventing the community from addressing the systemic disparities in female representation. However, a more recent figure of female OSS participation internationally calculated across 2005 to 2021 is 9.8%, with most being recent contributors, indicating that female participation may be growing.
Motivations
There are many motivations for contributing to the OSS community. For one, it is an opportunity to learn and practice multiple skills such as coding and other technology related abilities, but also fundamental skills such as communication and collaboration and practical skills needed to excel in technology related fields such as issue tracking or version control. Instead of learning through a classroom or a job, learning through contributing to OSS allows participants to learn at their own pace and follow what interests them. When contributing to OSS, the contributor can learn the current industry best practices, technology and trends and even have the opportunity to contribute to the next big innovation as OSS grows increasingly popular within the tech field. Contributing to OSS without payment means there is no threat of being fired, though reputations can take a hit. On the other hand, a huge motivation to contribute to OSS is the reputation gained as one grows one's public portfolio.
Disparities
Even though programming was originally seen as a female profession, there remains a large gap in computing. Social identity tends to be a large concern as women in the tech industry face insecurity about attracting unwanted male attention and harassment or being unfeminine in their technology knowledge, having a large impact on confidence. Some male tech participants make clear that they believe women fitting in within the culture is impossible, furthering the insecurity for women and their place in the tech industry. Additionally, even in a voluntary contribution environment like open source software, women tend to end up doing the less technical aspects of projects, such as manual testing or documentation despite women and men showing the same productivity in OSS contributions. Explicit biases include longer feedback time, more scrutinization of code and lower acceptance rate of code. Specifically in the open-source software community, women report that sexually offensive language is common and the women's identity as female is given more attention that as an OSS contributor Bias is hard to address due to the belief that gender should not matter, with most contributors feeling that women getting special treatment is unfair and success should be dependent on skill, preventing any changes to be more inclusive.
Adoption and application
Key projects
Open source software projects are built and maintained by a network of programmers, who may often be volunteers, and are widely used in free as well as commercial products.
Unix: Unix is an operating system created by AT&T that began as a precursor to open source software in that the free and open-source software revolution began when developers began trying to create operating systems without Unix code. Unix was created in the 1960s, before the commercialization of software and before the concept of open source software was necessary, therefore it was not considered a true open source software project. It started as a research project before being commercialized in the mid 1980s. Before its commercialization, it represented many of the ideals held by the Free and Open source software revolution, including the decentralized collaboration of global users, rolling releases and a community culture of distaste towards proprietary software.
BSD: Berkeley Software Distribution (BSD) is an operating system that began as a variant of Unix in 1978 that mixed Unix code with code from Berkeley labs to increase functionality. As BSD was focused on increasing functionality, it would publicly share its greatest innovations with the main Unix operating system. This is an example of the free public code sharing that is a central characteristic of FOSS today. As Unix became commercialized in the 1980s, developers or members of the community who did not support proprietary software began to focus on BSD and turning it into an operating system that did not include any of Unix's code. The final version of BSD was released in 1995.
GNU: GNU is a free operating system created by Richard Stallman in 1984 with its name meaning Gnu's Not Unix. The idea was to create a Unix alternative operating system that would be available for anyone to use and allow programmers to share code freely between them. However, the goal of GNU was not to only replace Unix, but to make a superior version that had more technological capabilities. It was released before the philosophical beliefs of the Free and Open source software revolution were truly defined. Because of its creation by prominent FOSS programmer Richard Stallman, GNU was heavily involved in FOSS activism, with one of the greatest achievements of GNU being the creation of the GNU General Public License or GPL, which allowed developers to release software that could be legally shared and modified.
Linux: Linux is an operating system kernel that was introduced in 1991 by Linus Torvalds. Linux was inspired by making a better version of the for profit operating service Minux. It was radically different than what other hackers were producing at the time due to it being totally free of cost and being decentralized. Later, Linux was put under the GPL license, allowing people to make money with Linux and bringing Linux into the FOSS community.
Apache: Apache began in 1995 as a collaboration between a group of developers releasing their own web server due to their frustration with NCSA HTTPd code base. The name Apache was used because of the several patches they applied to this code base. Within a year of its release, it became the worldwide leading web server. Soon, Apache came out with its own license, creating discord in the greater FOSS community, though ultimately proving successful. The Apache license allowed permitted members to directly access source code, a marked difference from GNU and Linux's approaches.
Extensions for non-software use
While the term open source applied originally only to the source code of software, it is now being applied to many other areas such as open-source ecology, a movement to decentralize technologies so that any human can use them. However, it is often misapplied to other areas that have different and competing principles, which overlap only partially.
The same principles that underlie open-source software can be found in many other ventures, such as open source, open content, and open collaboration.
This "culture" or ideology takes the view that the principles apply more generally to facilitate concurrent input of different agendas, approaches, and priorities, in contrast with more centralized models of development such as those typically used in commercial companies.
Value
More than 90 percent of companies use open-source software as a component of their proprietary software. The decision to use open-source software, or even engage with open-source projects to improve existing open-source software, is typically a pragmatic business decision. When proprietary software is in direct competition with an open-source alternative, research has found conflicting results on the effect of the competition on the proprietary product's price and quality.
For decades, some companies have made servicing of an open-source software product for enterprise users their business model. These companies control an open-source software product, and instead of charging for licensing or use, charge for improvements, integration, and other servicing. Software as a service (SaaS) products based on open-source components are increasingly common.
Open-source software is preferred for scientific applications, because it increases transparency and aids in the validation and acceptance of scientific results.
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277702 | https://en.wikipedia.org/wiki/Electron%20diffraction | Electron diffraction | Electron diffraction is a generic term for phenomena associated with changes in the direction of electron beams due to elastic interactions with atoms. It occurs due to elastic scattering, when there is no change in the energy of the electrons. The negatively charged electrons are scattered due to Coulomb forces when they interact with both the positively charged atomic core and the negatively charged electrons around the atoms. The resulting map of the directions of the electrons far from the sample is called a diffraction pattern, see for instance Figure 1. Beyond patterns showing the directions of electrons, electron diffraction also plays a major role in the contrast of images in electron microscopes.
This article provides an overview of electron diffraction and electron diffraction patterns, collective referred to by the generic name electron diffraction. This includes aspects of how in a general way electrons can act as waves, and diffract and interact with matter. It also involves the extensive history behind modern electron diffraction, how the combination of developments in the 19th century in understanding and controlling electrons in vacuum and the early 20th century developments with electron waves were combined with early instruments, giving birth to electron microscopy and diffraction in 1920–1935. While this was the birth, there have been a large number of further developments since then.
There are many types and techniques of electron diffraction. The most common approach is where the electrons transmit through a thin sample, from 1 nm to 100 nm (10 to 1000 atoms thick), where the results depending upon how the atoms are arranged in the material, for instance a single crystal, many crystals or different types of solids. Other cases such as larger repeats, no periodicity or disorder have their own characteristic patterns. There are many different ways of collecting diffraction information, from parallel illumination to a converging beam of electrons or where the beam is rotated or scanned across the sample which produce information that is often easier to interpret. There are also many other types of instruments. For instance, in a scanning electron microscope (SEM), electron backscatter diffraction can be used to determine crystal orientation across the sample. Electron diffraction patterns can also be used to characterize molecules using gas electron diffraction, liquids, surfaces using lower energy electrons, a technique called LEED, and by reflecting electrons off surfaces, a technique called RHEED.
There are also many levels of analysis of electron diffraction, including:
The simplest approximation using the de Broglie wavelength for electrons, where only the geometry is considered and often Bragg's law is invoked. This approach only considers the electrons far from the sample, a far-field or Fraunhofer approach.
The first level of more accuracy where it is approximated that the electrons are only scattered once, which is called kinematical diffraction and is also a far-field or Fraunhofer approach.
More complete and accurate explanations where multiple scattering is included, what is called dynamical diffraction (e.g. refs). These involve more general analyses using relativistically corrected Schrödinger equation methods, and track the electrons through the sample, being accurate both near and far from the sample (both Fresnel and Fraunhofer diffraction).
Electron diffraction is similar to x-ray and neutron diffraction. However, unlike x-ray and neutron diffraction where the simplest approximations are quite accurate, with electron diffraction this is not the case. Simple models give the geometry of the intensities in a diffraction pattern, but dynamical diffraction approaches are needed for accurate intensities and the positions of diffraction spots.
A primer on electron diffraction
All matter can be thought of as matter waves, from small particles such as electrons up to macroscopic objects – although it is impossible to measure any of the "wave-like" behavior of macroscopic objects. Waves can move around objects and create interference patterns, and a classic example is the Young's two-slit experiment shown in Figure 2, where a wave impinges upon two slits in the first of the two images (blue waves). After going through the slits there are directions where the wave is stronger, ones where it is weaker – the wave has been diffracted. If instead of two slits there are a number of small points then similar phenomena can occur as shown in the second image where the wave (red and blue) is coming in from the bottom right corner. This is comparable to diffraction of an electron wave where the small dots would be atoms in a small crystal, see also note. Note the strong dependence on the relative orientation of the crystal and the incoming wave.
Close to an aperture or atoms, often called the "sample", the electron wave would be described in terms of near field or Fresnel diffraction. This has relevance for imaging within electron microscopes, whereas electron diffraction patterns are measured far from the sample, which is described as far-field or Fraunhofer diffraction. A map of the directions of the electron waves leaving the sample will show high intensity (white) for favored directions, such as the three prominent ones in the Young's two-slit experiment of Figure 2, while the other directions will be low intensity (dark). Often there will be an array of spots (preferred directions) as in Figure 1 and the other figures shown later.
History
The historical background is divided into several subsections. The first is the general background to electrons in vacuum and the technological developments that led to cathode-ray tubes as well as vacuum tubes that dominated early television and electronics; the second is how these led to the development of electron microscopes; the last is work on the nature of electron beams and the fundamentals of how electrons behave, a key component of quantum mechanics and the explanation of electron diffraction.
Electrons in vacuum
Experiments involving electron beams occurred long before the discovery of the electron; ēlektron (ἤλεκτρον) is the Greek word for amber, which is connected to the recording of electrostatic charging by Thales of Miletus around 585 BCE, and possibly others even earlier.
In 1650, Otto von Guericke invented the vacuum pump allowing for the study of the effects of high voltage electricity passing through rarefied air. In 1838, Michael Faraday applied a high voltage between two metal electrodes at either end of a glass tube that had been partially evacuated of air, and noticed a strange light arc with its beginning at the cathode (negative electrode) and its end at the anode (positive electrode). Building on this, in the 1850s, Heinrich Geissler was able to achieve a pressure of around 10−3 atmospheres, inventing what became known as Geissler tubes. Using these tubes, while studying electrical conductivity in rarefied gases in 1859, Julius Plücker observed that the radiation emitted from the negatively charged cathode caused phosphorescent light to appear on the tube wall near it, and the region of the phosphorescent light could be moved by application of a magnetic field.
In 1869, Plücker's student Johann Wilhelm Hittorf found that a solid body placed between the cathode and the phosphorescence would cast a shadow on the tube wall, e.g. Figure 3. Hittorf inferred that there are straight rays emitted from the cathode and that the phosphorescence was caused by the rays striking the tube walls. In 1876 Eugen Goldstein showed that the rays were emitted perpendicular to the cathode surface, which differentiated them from the incandescent light. Eugen Goldstein dubbed them cathode rays. By the 1870s William Crookes and others were able to evacuate glass tubes below 10−6 atmospheres, and observed that the glow in the tube disappeared when the pressure was reduced but the glass behind the anode began to glow. Crookes was also able to show that the particles in the cathode rays were negatively charged and could be deflected by an electromagnetic field.
In 1897, Joseph Thomson measured the mass of these cathode rays, proving they were made of particles. These particles, however, were 1800 times lighter than the lightest particle known at that time – a hydrogen atom. These were originally called corpuscles and later named electrons by George Johnstone Stoney.
The control of electron beams that this work led to resulted in significant technology advances in electronic amplifiers and television displays.
Waves, diffraction and quantum mechanics
Independent of the developments for electrons in vacuum, at about the same time the components of quantum mechanics were being assembled. In 1924 Louis de Broglie in his PhD thesis Recherches sur la théorie des quanta introduced his theory of electron waves. He suggested that an electron around a nucleus could be thought of as standing waves, and that electrons and all matter could be considered as waves. He merged the idea of thinking about them as particles (or corpuscles), and of thinking of them as waves. He proposed that particles are bundles of waves (wave packets) that move with a group velocity and have an effective mass, see for instance Figure 4. Both of these depend upon the energy, which in turn connects to the wavevector and the relativistic formulation of Albert Einstein a few years before.
This rapidly became part of what was called by Erwin Schrödinger undulatory mechanics, now called the Schrödinger equation or wave mechanics. As stated by Louis de Broglie on September 8, 1927, in the preface to the German translation of his theses (in turn translated into English):M. Einstein from the beginning has supported my thesis, but it was M. E. Schrödinger who developed the propagation equations of a new theory and who in searching for its solutions has established what has become known as “Wave Mechanics”.
The Schrödinger equation combines the kinetic energy of waves and the potential energy due to, for electrons, the Coulomb potential. He was able to explain earlier work such as the quantization of the energy of electrons around atoms in the Bohr model, as well as many other phenomena. Electron waves as hypothesized by de Broglie were automatically part of the solutions to his equation, see also introduction to quantum mechanics and matter waves.
Both the wave nature and the undulatory mechanics approach were experimentally confirmed for electron beams by experiments from two groups performed independently, the first the Davisson–Germer experiment, the other by George Paget Thomson and Alexander Reid; see note for more discussion. Alexander Reid, who was Thomson's graduate student, performed the first experiments, but he died soon after in a motorcycle accident and is rarely mentioned. These experiments were rapidly followed by the first non-relativistic diffraction model for electrons by Hans Bethe based upon the Schrödinger equation, which is very close to how electron diffraction is now described. Significantly, Clinton Davisson and Lester Germer noticed that their results could not be interpreted using a Bragg's law approach as the positions were systematically different; the approach of Hans Bethe which includes the refraction due to the average potential yielded more accurate results. These advances in understanding of electron wave mechanics were important for many developments of electron-based analytical techniques such as Seishi Kikuchi's observations of lines due to combined elastic and inelastic scattering, gas electron diffraction developed by Herman Mark and Raymond Weil, diffraction in liquids by Louis Maxwell, and the first electron microscopes developed by Max Knoll and Ernst Ruska.
Electron microscopes and early electron diffraction
In order to have a practical microscope or diffractometer, just having an electron beam was not enough, it needed to be controlled. Many developments laid the groundwork of electron optics; see the paper by Chester J. Calbick for an overview of the early work. One significant step was the work of Heinrich Hertz in 1883 who made a cathode-ray tube with electrostatic and magnetic deflection, demonstrating manipulation of the direction of an electron beam. Others were focusing of electrons by an axial magnetic field by Emil Wiechert in 1899, improved oxide-coated cathodes which produced more electrons by Arthur Wehnelt in 1905 and the development of the electromagnetic lens in 1926 by Hans Busch.
Building an electron microscope involves combining these elements, similar to an optical microscope but with magnetic or electrostatic lenses instead of glass ones. To this day the issue of who invented the transmission electron microscope is controversial, as discussed by Thomas Mulvey and more recently by Yaping Tao. Extensive additional information can be found in the articles by Martin Freundlich, Reinhold Rüdenberg and Mulvey.
One effort was university based. In 1928, at the Technische Hochschule in Charlottenburg (now Technische Universität Berlin), (Professor of High Voltage Technology and Electrical Installations) appointed Max Knoll to lead a team of researchers to advance research on electron beams and cathode-ray oscilloscopes. The team consisted of several PhD students including Ernst Ruska. In 1931, Max Knoll and Ernst Ruska successfully generated magnified images of mesh grids placed over an anode aperture. The device, a replicate of which is shown in Figure 5, used two magnetic lenses to achieve higher magnifications, the first electron microscope. (Max Knoll died in 1969, so did not receive a share of the Nobel Prize in Physics in 1986.)
Apparently independent of this effort was work at Siemens-Schuckert by Reinhold Rudenberg. According to patent law (U.S. Patent No. 2058914 and 2070318, both filed in 1932), he is the inventor of the electron microscope, but it is not clear when he had a working instrument. He stated in a very brief article in 1932 that Siemens had been working on this for some years before the patents were filed in 1932, so his effort was parallel to the university effort. He died in 1961, so similar to Max Knoll, was not eligible for a share of the Nobel Prize.
These instruments could produce magnified images, but were not particularly useful for electron diffraction; indeed, the wave nature of electrons was not exploited during the development. Key for electron diffraction in microscopes was the advance in 1936 where showed that they could be used as micro-diffraction cameras with an aperture—the birth of selected area electron diffraction.
Less controversial was the development of LEED—the early experiments of Davisson and Germer used this approach. As early as 1929 Germer investigated gas adsorption, and in 1932 Harrison E. Farnsworth probed single crystals of copper and silver. However, the vacuum systems available at that time were not good enough to properly control the surfaces, and it took almost forty years before these became available. Similarly, it was not until about 1965 that Peter B. Sewell and M. Cohen demonstrated the power of RHEED in a system with a very well controlled vacuum.
Subsequent developments in methods and modelling
Despite early successes such as the determination of the positions of hydrogen atoms in NH4Cl crystals by W. E. Laschkarew and I. D. Usykin in 1933, boric acid by John M. Cowley in 1953 and orthoboric acid by William Houlder Zachariasen in 1954, electron diffraction for many years was a qualitative technique used to check samples within electron microscopes. John M Cowley explains in a 1968 paper: Thus was founded the belief, amounting in some cases almost to an article of faith, and persisting even to the present day, that it is impossible to interpret the intensities of electron diffraction patterns to gain structural information.This has changed, in transmission, reflection and for low energies. Some of the key developments (some of which are also described later) from the early days to 2023 have been:
Fast numerical methods based upon the Cowley–Moodie multislice algorithm, which only became possible once the fast Fourier transform (FFT) method was developed. With these and other numerical methods Fourier transforms are fast, and it became possible to calculate accurate, dynamical diffraction in seconds to minutes with laptops using widely available multislice programs.
Developments in the convergent-beam electron diffraction approach. Building on the original work of Walther Kossel and Gottfried Möllenstedt in 1939, it was extended by Peter Goodman and Gunter Lehmpfuhl, then mainly by the groups of John Steeds and Michiyoshi Tanaka who showed how to determine point groups and space groups. It can also be used for higher-level refinements of the electron density; for a brief history see CBED history. In many cases this is the best method to determine symmetry.
The development of new approaches to reduce dynamical effects such as precession electron diffraction and three-dimensional diffraction methods. Averaging over different directions has, empirically, been found to significantly reduce dynamical diffraction effects, e.g., see PED history for further details. Not only is it easier to identify known structures with this approach, it can also be used to solve unknown structures in some cases – see precession electron diffraction for further information.
The development of experimental methods exploiting ultra-high vacuum technologies (e.g. the approach described by in 1953) to better control surfaces, making LEED and RHEED more reliable and reproducible techniques. In the early days the surfaces were not well controlled; with these technologies they can both be cleaned and remain clean for hours to days, a key component of surface science.
Fast and accurate methods to calculate intensities for LEED so it could be used to determine atomic positions, for instance references. These have been extensively exploited to determine the structure of many surfaces, and the arrangement of foreign atoms on surfaces.
Methods to simulate the intensities in RHEED, so it can be used semi-quantitatively to understand surfaces during growth and thereby to control the resulting materials.
The development of advanced detectors for transmission electron microscopy such as charge-coupled device and direct electron detectors, which improve the accuracy and reliability of intensity measurements. These have efficiencies and accuracies that can be a thousand or more times that of the photographic film used in the earliest experiments, with the information available in real time rather than requiring photographic processing after the experiment.
Core elements of electron diffraction
Plane waves, wavevectors and reciprocal lattice
What is seen in an electron diffraction pattern depends upon the sample and also the energy of the electrons. The electrons need to be considered as waves, which involves describing the electron via a wavefunction, written in crystallographic notation (see notes and) as:for a position . This is a quantum mechanics description; one cannot use a classical approach. The vector is called the wavevector, has units of inverse nanometers, and the form above is called a plane wave as the term inside the exponential is constant on the surface of a plane. The vector is what is used when drawing ray diagrams, and in vacuum is parallel to the direction or, better, group velocity or probability current of the plane wave. For most cases the electrons are travelling at a respectable fraction of the speed of light, so rigorously need to be considered using relativistic quantum mechanics via the Dirac equation, which as spin does not normally matter can be reduced to the Klein–Gordon equation. Fortunately one can side-step many complications and use a non-relativistic approach based around the Schrödinger equation. Following Kunio Fujiwara and Archibald Howie, the relationship between the total energy of the electrons and the wavevector is written as:withwhere is the Planck constant, is a relativistic effective mass used to cancel out the relativistic terms for electrons of energy with the speed of light and the rest mass of the electron. The concept of effective mass occurs throughout physics (see for instance Ashcroft and Mermin), and comes up in the behavior of quasiparticles. A common one is the electron hole, which acts as if it is a particle with a positive charge and a mass similar to that of an electron, although it can be several times lighter or heavier. For electron diffraction the electrons behave as if they are non-relativistic particles of mass in terms of how they interact with the atoms.
The wavelength of the electrons in vacuum is from the above equationsand can range from about , roughly the size of an atom, down to a thousandth of that. Typically the energy of the electrons is written in electronvolts (eV), the voltage used to accelerate the electrons; the actual energy of each electron is this voltage times the electron charge. For context, the typical energy of a chemical bond is a few eV; electron diffraction involves electrons up to .
The magnitude of the interaction of the electrons with a material scales asWhile the wavevector increases as the energy increases, the change in the effective mass compensates this so even at the very high energies used in electron diffraction there are still significant interactions.
The high-energy electrons interact with the Coulomb potential, which for a crystal can be considered in terms of a Fourier series (see for instance Ashcroft and Mermin), that iswith a reciprocal lattice vector and the corresponding Fourier coefficient of the potential. The reciprocal lattice vector is often referred to in terms of Miller indices , a sum of the individual reciprocal lattice vectors with integers in the form:(Sometimes reciprocal lattice vectors are written as , , and see note.) The contribution from the needs to be combined with what is called the shape function (e.g.), which is the Fourier transform of the shape of the object. If, for instance, the object is small in one dimension then the shape function extends far in that direction in the Fourier transform—a reciprocal relationship.
Around each reciprocal lattice point one has this shape function. How much intensity there will be in the diffraction pattern depends upon the intersection of the Ewald sphere, that is energy conservation, and the shape function around each reciprocal lattice point—see Figure 6, 20 and 22. The vector from a reciprocal lattice point to the Ewald sphere is called the excitation error .
For transmission electron diffraction the samples used are thin, so most of the shape function is along the direction of the electron beam. For both LEED and RHEED the shape function is mainly normal to the surface of the sample. In LEED this results in (a simplification) back-reflection of the electrons leading to spots, see Figure 20 and 21 later, whereas in RHEED the electrons reflect off the surface at a small angle and typically yield diffraction patterns with streaks, see Figure 22 and 23 later. By comparison, with both x-ray and neutron diffraction the scattering is significantly weaker, so typically requires much larger crystals, in which case the shape function shrinks to just around the reciprocal lattice points, leading to simpler Bragg's law diffraction.
For all cases, when the reciprocal lattice points are close to the Ewald sphere (the excitation error is small) the intensity tends to be higher; when they are far away it tends to be smaller. The set of diffraction spots at right angles to the direction of the incident beam are called the zero-order Laue zone (ZOLZ) spots, as shown in Figure 6. One can also have intensities further out from reciprocal lattice points which are in a higher layer. The first of these is called the first order Laue zone (FOLZ); the series is called by the generic name higher order Laue zone (HOLZ).
The result is that the electron wave after it has been diffracted can be written as an integral over different plane waves:that is a sum of plane waves going in different directions, each with a complex amplitude . (This is a three dimensional integral, which is often written as rather than .) For a crystalline sample these wavevectors have to be of the same magnitude for elastic scattering (no change in energy), and are related to the incident direction by (see Figure 6)
A diffraction pattern detects the intensitiesFor a crystal these will be near the reciprocal lattice points typically forming a two dimensional grid. Different samples and modes of diffraction give different results, as do different approximations for the amplitudes .
A typical electron diffraction pattern in TEM and LEED is a grid of high intensity spots (white) on a dark background, approximating a projection of the reciprocal lattice vectors, see Figure 1, 9, 10, 11, 14 and 21 later. There are also cases which will be mentioned later where diffraction patterns are not periodic, see Figure 15, have additional diffuse structure as in Figure 16, or have rings as in Figure 12, 13 and 24. With conical illumination as in CBED they can also be a grid of discs, see Figure 7, 9 and 18. RHEED is slightly different, see Figure 22, 23. If the excitation errors were zero for every reciprocal lattice vector, this grid would be at exactly the spacings of the reciprocal lattice vectors. This would be equivalent to a Bragg's law condition for all of them. In TEM the wavelength is small and this is close to correct, but not exact. In practice the deviation of the positions from a simple Bragg's law interpretation is often neglected, particularly if a column approximation is made (see below).
Kinematical diffraction
In Kinematical theory an approximation is made that the electrons are only scattered once. For transmission electron diffraction it is common to assume a constant thickness , and also what is called the Column Approximation (e.g. references and further reading). For a perfect crystal the intensity for each diffraction spot is then:where is the magnitude of the excitation error along z, the distance along the beam direction (z-axis by convention) from the diffraction spot to the Ewald sphere, and is the structure factor:the sum being over all the atoms in the unit cell with the form factors, the reciprocal lattice vector, is a simplified form of the Debye–Waller factor, and is the wavevector for the diffraction beam which is:for an incident wavevector of , as in Figure 6 and above. The excitation error comes in as the outgoing wavevector has to have the same modulus (i.e. energy) as the incoming wavevector . The intensity in transmission electron diffraction oscillates as a function of thickness, which can be confusing; there can similarly be intensity changes due to variations in orientation and also structural defects such as dislocations. If a diffraction spot is strong it could be because it has a larger structure factor, or it could be because the combination of thickness and excitation error is "right". Similarly the observed intensity can be small, even though the structure factor is large. This can complicate interpretation of the intensities. By comparison, these effects are much smaller in x-ray diffraction or neutron diffraction because they interact with matter far less and often Bragg's law is adequate.
This form is a reasonable first approximation which is qualitatively correct in many cases, but more accurate forms including multiple scattering (dynamical diffraction) of the electrons are needed to properly understand the intensities.
Dynamical diffraction
While kinematical diffraction is adequate to understand the geometry of the diffraction spots, it does not correctly give the intensities and has a number of other limitations. For a more complete approach one has to include multiple scattering of the electrons using methods that date back to the early work of Hans Bethe in 1928. These are based around solutions of the Schrödinger equation using the relativistic effective mass described earlier. Even at very high energies dynamical diffraction is needed as the relativistic mass and wavelength partially cancel, so the role of the potential is larger than might be thought.
The main components of current dynamical diffraction of electrons include:
Taking into account the scattering back into the incident beam both from diffracted beams and between all others, not just single scattering from the incident beam to diffracted beams. This is important even for samples which are only a few atoms thick.
Modelling at least semi-empirically the role of inelastic scattering by an imaginary component of the potential, also called an "optical potential". There is always inelastic scattering, and often it can have a major effect on both the background and sometimes the details, see Figure 7 and 18.
Higher-order numerical approaches to calculate the intensities such as multislice, matrix methods which are called Bloch-wave approaches or muffin-tin approaches. With these diffraction spots which are not present in kinematical theory can be present, e.g.
Contributions to the diffraction from elastic strain and crystallographic defects, and also what Jens Lindhard called the string potential.
For transmission electron microscopes effects due to variations in the thickness of the sample and the normal to the surface.
Both in the geometry of scattering and calculations, for both LEED and RHEED, effects due to the presence of surface steps, surface reconstructions and other atoms at the surface. Often these change the diffraction details significantly.
For LEED, use more careful analyses of the potential because contributions from exchange terms can be important. Without these the calculations may not be accurate enough.
Kikuchi lines
Kikuchi lines, first observed by Seishi Kikuchi in 1928, are linear features created by electrons scattered both inelastically and elastically. As the electron beam interacts with matter, the electrons are diffracted via elastic scattering, and also scattered inelastically losing part of their energy. These occur simultaneously, and cannot be separated – according to the Copenhagen interpretation of quantum mechanics, only the probabilities of electrons at detectors can be measured. These electrons form Kikuchi lines which provide information on the orientation.
Kikuchi lines come in pairs forming Kikuchi bands, and are indexed in terms of the crystallographic planes they are connected to, with the angular width of the band equal to the magnitude of the corresponding diffraction vector . The position of Kikuchi bands is fixed with respect to each other and the orientation of the sample, but not against the diffraction spots or the direction of the incident electron beam. As the crystal is tilted, the bands move on the diffraction pattern. Since the position of Kikuchi bands is quite sensitive to crystal orientation, they can be used to fine-tune a zone-axis orientation or determine crystal orientation. They can also be used for navigation when changing the orientation between zone axes connected by some band, an example of such a map produced by combining many local sets of experimental Kikuchi patterns is in Figure 8; Kikuchi maps are available for many materials.
Types and techniques
In a transmission electron microscope
Electron diffraction in a TEM exploits controlled electron beams using electron optics. Different types of diffraction experiments, for instance Figure 9, provide information such as lattice constants, symmetries, and sometimes to solve an unknown crystal structure.
It is common to combine it with other methods, for instance images using selected diffraction beams, high-resolution images showing the atomic structure, chemical analysis through energy-dispersive x-ray spectroscopy, investigations of electronic structure and bonding through electron energy loss spectroscopy, and studies of the electrostatic potential through electron holography; this list is not exhaustive. Compared to x-ray crystallography, TEM analysis is significantly more localized and can be used to obtain information from tens of thousands of atoms to just a few or even single atoms.
Formation of a diffraction pattern
In TEM, the electron beam passes through a thin film of the material as illustrated in Figure 10. Before and after the sample the beam is manipulated by the electron optics including magnetic lenses, deflectors and apertures; these act on the electrons similar to how glass lenses focus and control light. Optical elements above the sample are used to control the incident beam which can range from a wide and parallel beam to one which is a converging cone and can be smaller than an atom, 0.1 nm. As it interacts with the sample, part of the beam is diffracted and part is transmitted without changing its direction. This occurs simultaneously as electrons are everywhere until they are detected (wavefunction collapse) according to the Copenhagen interpretation.
Below the sample, the beam is controlled by another set of magnetic lneses and apertures. Each set of initially parallel rays (a plane wave) is focused by the first lens (objective) to a point in the back focal plane of this lens, forming a spot on a detector; a map of these directions, often an array of spots, is the diffraction pattern. Alternatively the lenses can form a magnified image of the sample. Herein the focus is on collecting a diffraction pattern; for other information see the pages on TEM and scanning transmission electron microscopy.
Selected area electron diffraction
The simplest diffraction technique in TEM is selected area electron diffraction (SAED) where the incident beam is wide and close to parallel. An aperture is used to select a particular region of interest from which the diffraction is collected. These apertures are part of a thin foil of a heavy metal such as tungsten which has a number of small holes in it. This way diffraction information can be limited to, for instance, individual crystallites. Unfortunately the method is limited by the spherical aberration of the objective lens, so is only accurate for large grains with tens of thousands of atoms or more; for smaller regions a focused probe is needed.
If a parallel beam is used to acquire a diffraction pattern from a single-crystal, the result is similar to a two-dimensional projection of the crystal reciprocal lattice. From this one can determine interplanar distances and angles and in some cases crystal symmetry, particularly when the electron beam is down a major zone axis, see for instance the database by Jean-Paul Morniroli. However, projector lens aberrations such as barrel distortion as well as dynamical diffraction effects (e.g.) cannot be ignored. For instance, certain diffraction spots which are not present in x-ray diffraction can appear, for instance those due to Gjønnes-Moodie extinction conditions.
If the sample is tilted relative to the electron beam, different sets of crystallographic planes contribute to the pattern yielding different types of diffraction patterns, approximately different projections of the reciprocal lattice, see Figure 11. This can be used to determine the crystal orientation, which in turn can be used to set the orientation needed for a particular experiment. Furthermore, a series of diffraction patterns varying in tilt can be acquired and processed using a diffraction tomography approach. There are ways to combine this with direct methods algorithms using electrons and other methods such as charge flipping, or automated diffraction tomography to solve crystal structures.
Polycrystalline pattern
Diffraction patterns depend on whether the beam is diffracted by one single crystal or by a number of differently oriented crystallites, for instance in a polycrystalline material. If there are many contributing crystallites, the diffraction image is a superposition of individual crystal patterns, see Figure 12. With a large number of grains this superposition yields diffraction spots of all possible reciprocal lattice vectors. This results in a pattern of concentric rings as shown in Figure 12 and 13.
Textured materials yield a non-uniform distribution of intensity around the ring, which can be used to discriminate between nanocrystalline and amorphous phases. However, diffraction often cannot differentiate between very small grain polycrystalline materials and truly random order amorphous. Here high-resolution transmission electron microscopy and fluctuation electron microscopy can be more powerful, although this is still a topic of continuing development.
Multiple materials and double diffraction
In simple cases there is only one grain or one type of material in the area used for collecting a diffraction pattern. However, often there is more than one. If they are in different areas then the diffraction pattern will be a combination. In addition there can be one grain on top of another, in which case the electrons that go through the first are diffracted by the second. Electrons have no memory (like many of us), so after they have gone through the first grain and been diffracted, they traverse the second as if their current direction was that of the incident beam. This leads to diffraction spots which are the vector sum of those of the two (or even more) reciprocal lattices of the crystals, and can lead to complicated results. It can be difficult to know if this is real and due to some novel material, or just a case where multiple crystals and diffraction is leading to odd results.
Bulk and surface superstructures
Many materials have relatively simple structures based upon small unit cell vectors (see also note). There are many others where the repeat is some larger multiple of the smaller unit cell (subcell) along one or more direction, for instance . which has larger dimensions in two directions. These superstructures can arise from many reasons:
Larger unit cells due to electronic ordering which leads to small displacements of the atoms in the subcell. One example is antiferroelectricity ordering.
Chemical ordering, that is different atom types at different locations of the subcell.
Magnetic order of the spins. These may be in opposite directions on some atoms, leading to what is called antiferromagnetism.
In addition to those which occur in the bulk, superstructures can also occur at surfaces. When half the material is (nominally) removed to create a surface, some of the atoms will be under coordinated. To reduce their energy they can rearrange. Sometimes these rearrangements are relatively small; sometimes they are quite large. Similar to a bulk superstructure there will be additional, weaker diffraction spots. One example is for the silicon (111) surface, where there is a supercell which is seven times larger than the simple bulk cell in two directions. This leads to diffraction patterns with additional spots some of which are marked in Figure 14. Here the (220) are stronger bulk diffraction spots, and the weaker ones due to the surface reconstruction are marked 7 × 7—see note for convention comments.
Aperiodic materials
In an aperiodic crystal the structure can no longer be simply described by three different vectors in real or reciprocal space. In general there is a substructure describable by three (e.g. ), similar to supercells above, but in addition there is some additional periodicity (one to three) which cannot be described as a multiple of the three; it is a genuine additional periodicity which is an irrational number relative to the subcell lattice. The diffraction pattern can then only be described by more than three indices.
An extreme example of this is for quasicrystals, which can be described similarly by a higher number of Miller indices in reciprocal space—but not by any translational symmetry in real space. An example of this is shown in Figure 15 for an Al–Cu–Fe–Cr decagonal quasicrystal grown by magnetron sputtering on a sodium chloride substrate and then lifted off by dissolving the substrate with water. In the pattern there are pentagons which are a characteristic of the aperiodic nature of these materials.
Diffuse scattering
A further step beyond superstructures and aperiodic materials is what is called diffuse scattering in electron diffraction patterns due to disorder, which is also known for x-ray or neutron scattering. This can occur from inelastic processes, for instance, in bulk silicon the atomic vibrations (phonons) are more prevalent along specific directions, which leads to streaks in diffraction patterns. Sometimes it is due to arrangements of point defects. Completely disordered substitutional point defects lead to a general background which is called Laue monotonic scattering. Often there is a probability distribution for the distances between point defects or what type of substitutional atom there is, which leads to distinct three-dimensional intensity features in diffraction patterns. An example of this is for a Nb0.83CoSb sample, with the diffraction pattern shown in Figure 16. Because of the vacancies at the niobium sites, there is diffuse intensity with snake-like structure due to correlations of the distances between vacancies and also the relaxation of Co and Sb atoms around these vacancies.
Convergent beam electron diffraction
In convergent beam electron diffraction (CBED), the incident electrons are normally focused in a converging cone-shaped beam with a crossover located at the sample, e.g. Figure 17, although other methods exist. Unlike the parallel beam, the convergent beam is able to carry information from the sample volume, not just a two-dimensional projection available in SAED. With convergent beam there is also no need for the selected area aperture, as it is inherently site-selective since the beam crossover is positioned at the object plane where the sample is located.
A CBED pattern consists of disks arranged similar to the spots in SAED. Intensity within the disks represents dynamical diffraction effects and symmetries of the sample structure, see Figure 7 and 18. Even though the zone axis and lattice parameter analysis based on disk positions does not significantly differ from SAED, the analysis of disks content is more complex and simulations based on dynamical diffraction theory is often required. As illustrated in Figure 18, the details within the disk change with sample thickness, as does the inelastic background. With appropriate analysis CBED patterns can be used for indexation of the crystal point group, space group identification, measurement of lattice parameters, thickness or strain.
The disk diameter can be controlled using the microscope optics and apertures. The larger is the angle, the broader the disks are with more features. If the angle is increased to significantly, the disks begin to overlap. This is avoided in large angle convergent electron beam diffraction (LACBED) where the sample is moved upwards or downwards. There are applications, however, where the overlapping disks are beneficial, for instance with a ronchigram. It is a CBED pattern, often but not always of an amorphous material, with many intentionally overlapping disks providing information about the optical aberrations of the electron optical system.
Precession electron diffraction
Precession electron diffraction (PED), invented by Roger Vincent and Paul Midgley in 1994, is a method to collect electron diffraction patterns in a transmission electron microscope (TEM). The technique involves rotating (precessing) a tilted incident electron beam around the central axis of the microscope, compensating for the tilt after the sample so a spot diffraction pattern is formed, similar to a SAED pattern. However, a PED pattern is an integration over a collection of diffraction conditions, see Figure 19. This integration produces a quasi-kinematical diffraction pattern that is more suitable as input into direct methods algorithms using electrons to determine the crystal structure of the sample. Because it avoids many dynamical effects it can also be used to better identify crystallographic phases.
4D STEM
4D scanning transmission electron microscopy (4D STEM) is a subset of scanning transmission electron microscopy (STEM) methods which uses a pixelated electron detector to capture a convergent beam electron diffraction (CBED) pattern at each scan location; see the main page for further information. This technique captures a 2 dimensional reciprocal space image associated with each scan point as the beam rasters across a 2 dimensional region in real space, hence the name 4D STEM. Its development was enabled by better STEM detectors and improvements in computational power. The technique has applications in diffraction contrast imaging, phase orientation and identification, strain mapping, and atomic resolution imaging among others; it has become very popular and rapidly evolving from about 2020 onwards.
The name 4D STEM is common in literature, however it is known by other names: 4D STEM EELS, ND STEM (N- since the number of dimensions could be higher than 4), position resolved diffraction (PRD), spatial resolved diffractometry, momentum-resolved STEM, "nanobeam precision electron diffraction", scanning electron nano diffraction, nanobeam electron diffraction, or pixelated STEM. Most of these are the same, although there are instances such as momentum-resolved STEM where the emphasis can be very different.
Low-energy electron diffraction (LEED)
Low-energy electron diffraction (LEED) is a technique for the determination of the surface structure of single-crystalline materials by bombardment with a collimated beam of low-energy electrons (30–200 eV). In this case the Ewald sphere leads to approximately back-reflection, as illustrated in Figure 20, and diffracted electrons as spots on a fluorescent screen as shown in Figure 21; see the main page for more information and references. It has been used to solve a very large number of relatively simple surface structures of metals and semiconductors, plus cases with simple chemisorbants. For more complex cases transmission electron diffraction or surface x-ray diffraction have been used, often combined with scanning tunneling microscopy and density functional theory calculations.
LEED may be used in one of two ways:
Qualitatively, where the diffraction pattern is recorded and analysis of the spot positions gives information on the symmetry of the surface structure. In the presence of an adsorbate the qualitative analysis may reveal information about the size and rotational alignment of the adsorbate unit cell with respect to the substrate unit cell.
Quantitatively, where the intensities of diffracted beams are recorded as a function of incident electron beam energy to generate the so-called I–V curves. By comparison with theoretical curves, these may provide accurate information on atomic positions on the surface.
Reflection high-energy electron diffraction (RHEED)
Reflection high energy electron diffraction (RHEED), is a technique used to characterize the surface of crystalline materials by reflecting electrons off a surface. As illustrated for the Ewald sphere construction in Figure 22, it uses mainly the higher-order Laue zones which have a reflection component. An experimental diffraction pattern is shown in Figure 23 and shows both rings from the higher-order Laue zones and streaky spots. RHEED systems gather information only from the surface layers of the sample, which distinguishes RHEED from other materials characterization methods that also rely on diffraction of electrons. Transmission electron microscopy samples mainly the bulk of the sample, although in special cases it can provide surface information. Low-energy electron diffraction (LEED) is also surface sensitive, and achieves surface sensitivity through the use of low energy electrons. The main uses of RHEED to date have been during thin film growth, as the geometry is amenable to simultaneous collection of the diffraction data and deposition. It can, for instance, be used to monitor surface roughness during growth by looking at both the shapes of the streaks in the diffraction pattern as well as variations in the intensities.
Gas electron diffraction
Gas electron diffraction (GED) can be used to determine the geometry of molecules in gases. A gas carrying the molecules is exposed to the electron beam, which is diffracted by the molecules. Since the molecules are randomly oriented, the resulting diffraction pattern consists of broad concentric rings, see Figure 24. The diffraction intensity is a sum of several components such as background, atomic intensity or molecular intensity.
In GED the diffraction intensities at a particular diffraction angle is described via a scattering variable defined asThe total intensity is then given as a sum of partial contributions:where results from scattering by individual atoms, by pairs of atoms and by atom triplets. Intensity corresponds to the background which, unlike the previous contributions, must be determined experimentally. The intensity of atomic scattering is defined aswhere , is the distance between the scattering object detector, is the intensity of the primary electron beam and is the scattering amplitude of the atom of the molecular structure in the experiment. is the main contribution and easily obtained for known gas composition. Note that the vector used here is not the same as the excitation error used in other areas of diffraction, see earlier.
The most valuable information is carried by the intensity of molecular scattering , as it contains information about the distance between all pairs of atoms in the molecule. It is given bywhere is the distance between two atoms, is the mean square amplitude of vibration between the two atoms, similar to a Debye–Waller factor, is the anharmonicity constant and a phase factor which is important for atomic pairs with very different nuclear charges. The summation is performed over all atom pairs. Atomic triplet intensity is negligible in most cases. If the molecular intensity is extracted from an experimental pattern by subtracting other contributions, it can be used to match and refine a structural model against the experimental data.
Similar methods of analysis have also been applied to analyze electron diffraction data from liquids.
In a scanning electron microscope
In a scanning electron microscope the region near the surface can be mapped using an electron beam that is scanned in a grid across the sample. A diffraction pattern can be recorded using electron backscatter diffraction (EBSD), as illustrated in Figure 25, captured with a camera inside the microscope. A depth from a few nanometers to a few microns, depending upon the electron energy used, is penetrated by the electrons, some of which are diffracted backwards and out of the sample. As result of combined inelastic and elastic scattering, typical features in an EBSD image are Kikuchi lines. Since the position of Kikuchi bands is highly sensitive to the crystal orientation, EBSD data can be used to determine the crystal orientation at particular locations of the sample. The data are processed by software yielding two-dimensional orientation maps. As the Kikuchi lines carry information about the interplanar angles and distances and, therefore, about the crystal structure, they can also be used for phase identification or strain analysis.
| Physical sciences | Quantum mechanics | Physics |
277709 | https://en.wikipedia.org/wiki/Vitellaria | Vitellaria | Vitellaria paradoxa (formerly Butyrospermum parkii), commonly known as shea tree, shi tree (, also ), or vitellaria, is a tree of the family Sapotaceae. It is the only species in the genus Vitellaria, and is indigenous to Africa.
The shea fruit consists of a thin, tart, nutritious pulp that surrounds a relatively large, oil-rich seed from which shea butter is extracted. It is a deciduous tree usually tall, but has reached and a trunk diameter of .
The shea tree is a traditional African food plant. It has been claimed to have potential to improve nutrition, boost food supply in the "annual hungry season", foster rural development, and support sustainable land care.
Description
The tree starts bearing its first fruit when it is 10 to 15 years old; full production is attained when the tree is about 20 to 30 years old. It then produces nuts for up to 200 years.
The fruits resemble large plums 4 to 8 centimetres long weighing between 10 and 57 grams each. These fruits take 4 to 6 months to ripen; the average yield is of fresh fruit per tree, with optimum yields up to . Each kilogram of fruit gives approximately of dry seeds. The fruit is edible.
Nomenclature and taxonomy
Vitellaria is a monotypic genus, i.e., it comprises only one species. The species has variously been known botanically as Vitellaria paradoxa, Butyrospermum parkii, and Butyrospermum paradoxum. Many botanical works from the late 19th and much of the 20th centuries used the name Butyrospermum parkii, which is still commonly found in the cosmetics trade. However, Vitellaria paradoxa is the oldest name (published in 1807) and has been generally used in recent decades, as necessitated by the rules of botanical nomenclature; efforts in 1962 to make Butyrospermum the official scientific name for the genus (i.e., to "conserve" the name) were unsuccessful.
The species comprises two subspecies:
Vitellaria paradoxa subsp. paradoxa (roughly from the Nigeria-Cameroon border westward).
Vitellaria paradoxa subsp. nilotica (Kotschy) A.N. Henry & Chithra & N.C. Nair (roughly from the Nigeria-Cameroon border eastward).
Distribution and habitat
The shea tree grows naturally in the wild in the dry savannah belt of West and South from Senegal in the west to Sudan and South Sudan in the east, and onto the foothills of the Ethiopian highlands. It occurs in 19 countries across the African continent, namely Benin, Burkina Faso, Cameroon, Central African Republic, Chad, Ethiopia, Ghana, Guinea Bissau, Ivory Coast, Mali, Niger, Nigeria, Senegal, Sierra Leone, South Sudan, Sudan, Togo, Uganda, Democratic Republic of the Congo, and Guinea. The habitat area extends over more than .
A testa found at the site of the medieval village of Saouga is evidence of shea butter production by the 14th century.
Uses
Shea butter has many uses and may or may not be refined. In the West it is most commonly used as an emollient in cosmetics and is less commonly used in food. Throughout Africa it is used extensively for food, is a major source of dietary fat, and for medicinal purposes. In Ghana and Nigeria, shea butter is a major ingredient for making the African black soap.
The edible protein-rich caterpillars of the moth Cirina butyrospermi which feed solely on its leaves are widely collected and eaten raw, dried or fried.
Composition of shea butter
Shea butter extract is a complex fat that in addition to many nonsaponifiable components (substances that cannot be fully converted into soap by treatment with alkali) contains the following fatty acids: oleic acid (40–60%), stearic acid (20–50%), linoleic acid (3–11%), palmitic acid (2–9%), linolenic acid (<1%) and arachidic acid (<1%). It also contains the vitamins A, E and F.
Etymology
The common name is shíyiri (in N'Ko: ) or shísu (, lit. "shea tree") in the Bambara language of Mali. This is the origin of the English word, whose primary pronunciation is (rhyming with "tea"), although the pronunciation (rhyming with "day") is common and is listed second in major dictionaries. The tree is called ghariti in the Wolof language of Senegal, which is the origin of the French name of the tree and the butter, karité.
In Hausa language the tree is called Kaɗe or Kaɗanya. Indeed, the shea tree is so indispensable in Mole-Dagbang culinary and ethno-botanical practices that the Northern Ghanaian city of Tamale etymologically derives its name from the more traditional Dagomba name 'Tama-yile' (meaning 'Home of Shea nuts').
The tree was formerly classified in the genus Butyrospermum, meaning "butter seed". The species name parkii honors Scottish explorer Mungo Park, who learned of the tree while exploring Senegal. Park's Scottish origin is reflected in the English word shea, with a final -ea.
| Biology and health sciences | Nuts | Plants |
277833 | https://en.wikipedia.org/wiki/Phoneutria | Phoneutria | Phoneutria is a genus of spiders in the family Ctenidae. They are mainly found in northern South America, with one species in Central America. Members of the genus are commonly referred to as Brazilian wandering spiders. Other English names include armed spiders (armadeiras in Brazilian Portuguese) and banana spiders (a name shared with several others).
Description
The spiders in the genus can grow to have a leg span of . Their body length ranges from . While some other araneomorph spiders have a longer leg span, the largest Phoneutria species have the longest body and the greatest body weight in this group. The genus is distinguished from other related genera such as Ctenus by the presence of dense prolateral scopulae (a dense brush of fine hairs) on the pedipalp tibiae and tarsi in both sexes. Phoneutria are easily confused with several other non-medically significant ctenids, especially Cupiennius, in which the recently described C. chiapanensis also has bright red hairs on the chelicerae. Additionally, some Phoneutria species lack red hairs on the chelicerae, making it an unreliable identification feature. The presence of a dark linear stripe or stripes on the frontal (dorsal) palps and presence of a single thin black line running anterior-posterior along the dorsal carapace may help identify Phoneutria. Other features are the strong ventral marking on the underside of the legs with contrasting dark mid-segments and lighter joints, and the pattern on the ventral (underside) of the abdomen with several rows of black dots, or an overall reddish colour.
The characteristic defensive posture with frontal legs held high is an especially good indicator to confirm a specimen is Phoneutria, especially alongside correct colour patterns. During the defensive display the body is lifted up into an erect position, the first two pairs of legs are lifted high (revealing the conspicuous black/light-banded pattern on the leg underside), while the spider sways from side to side with hind legs in a cocked position.
Taxonomy
The genus Phoneutria was started by Maximilian Perty in 1833. The genus name is from the Greek φονεύτρια, meaning "murderess". Perty placed two species in the genus: Phoneutria rufibarbis and Phoneutria fera. The former is treated as a nomen dubium; the latter is the type species of the genus.
Species
, the World Spider Catalog accepted the following species:
Phoneutria bahiensis Simó & Brescovit, 2001 – Brazil
Phoneutria boliviensis (F. O. Pickard-Cambridge, 1897) – Central, South America
Phoneutria depilata (Strand, 1909) – Guatemala, Honduras, Nicaragua, Costa Rica, Panama, Colombia, Ecuador
Phoneutria eickstedtae Martins & Bertani, 2007 – Brazil
Phoneutria fera Perty, 1833 (type) – Colombia, Ecuador, Peru, Brazil, Suriname, Guyana
Phoneutria keyserlingi (F. O. Pickard-Cambridge, 1897) – Brazil
Phoneutria nigriventer (Keyserling, 1891) – Brazil, Uruguay, Paraguay, Argentina
Phoneutria pertyi (F. O. Pickard-Cambridge, 1897) – Brazil
Phoneutria reidyi (F. O. Pickard-Cambridge, 1897) – Colombia, Venezuela, Peru, Brazil, Guyana
Behaviour
Wandering spiders are so-called because they wander the jungle floor at night, rather than residing in a lair or maintaining a web. During the day they hide inside termite mounds, under fallen logs and rocks, and in banana plants (hence the "banana spider" nickname) and bromeliads. P. nigriventer is known to hide in dark and moist places in or near human dwellings.
P. nigriventer mates during the dry season from April to June, which leads to frequent observations of the species during this time.
Distribution
Phoneutria are found in forests from Costa Rica southwards throughout South America east of the Andes including Colombia, Venezuela, the Guianas, Ecuador, Peru, Bolivia, Brazil, Paraguay, and into northern Argentina. Three species (P. reidyi, P. boliviensis and P. fera) are found in the Amazon region, one species (P. fera) is restricted to the Amazon, and one (P. depilata) ranges into Central America in Panama and Costa Rica. The remaining species are restricted to Atlantic Forest of Argentina, Paraguay and Brazil, including forest fragments in the Cerrado savanna. In Brazil, Phoneutria is only absent in the northeastern region north of Salvador, Bahia.
Phoneutria has been introduced to Chile and Uruguay.
Banana shipments
These spiders acquired their other common name, "banana spider", because it is claimed that they are occasionally found in shipments of bananas, though the number of reports is exaggerated due to common misidentifications of unrelated spiders. A survey of spiders found in international shipments to North America revealed that only 7 of 135 spiders were Phoneutria species, six being Phoneutria boliviensis from bananas and one Phoneutria nigriventer from a shipment of electrical parts. Spiders from genera such as Cupiennius had been misidentified by experienced arachnologists. Cases continue to be reported but without evidence of expert identification. In 2005, a man was bitten in Bridgwater, England by a spider in a shipment of bananas and, in 2014, a south London family photographed a spider that they claim was in a bunch of bananas delivered to their home.
Medical significance
The genus Phoneutria includes some of the relatively few species of spiders known to present a threat to humans. Danger to humans does not only comprise toxicity, but also factors such as the spider's capacity to deliver a sufficient dose of venom, a disposition that makes a bite likely, and proximity to human habitation. These spiders' wandering nature is another reason they are considered so dangerous. In densely populated areas, Phoneutria species usually search for cover and dark places to hide during daytime, leading them to hide in houses, clothes, cars, boots, boxes and log piles, where they may bite if accidentally disturbed.
Spider mouthparts are adapted to envenomate very small prey; they are not well-adapted to attacking large mammals such as humans. Some believe that various spiders like Phoneutria, that use venom mainly to kill prey, can deliver a "dry" bite in defense to purposely conserve their venom, as opposed to a more primitive spider like Atrax that usually delivers a full load. A study in March 2009 suggests that Phoneutria inject venom in approximately one-third of their bites, and only a small quantity in one-third of those cases. Another study similarly suggested that only 2.3% of bites (mainly in children) were serious enough to require antivenom. Research in this area is hindered by the difficulty of identifying particular species. Nevertheless, there are a few well-attested instances of death. In one case, a single spider killed two children in São Sebastião. The spider was positively identified as a Phoneutria by Wolfgang Bücherl. Fatalities are usually attributed to respiratory arrest, secondary to systemic effects, or directly to envenoming. Systemic effects occur in 9% to 27% of cases; symptoms at the sites are more frequent, occurring from 83% to 96% of cases. The severity of the cases can be related to the sex of the spider, since the male produces less venom and is less lethal than the females, except for P. boliviensis, where the male is more toxic. Symptoms may appear within 10 to 20 minutes after the bite, and death within two to six hours, where severe pain radiates to the rest of the limb, systemic effects include tachycardia, increased blood pressure, vertigo, fever, sweating, visual disturbances, nausea, vomiting, difficulty breathing and paralysis. Death is usually caused by respiratory arrest. These spiders seem to produce a smaller amount of venom during cold months (June to September) a minimum amount of 0.03 mg, an average of 0.44 mg and a maximum of 1.84 mg, during the summer months. The maximum amount among individuals was 3.10 mg (October 26), 4 mg (November 3), 5 mg (November 4) and 8 mg (October 31); 7 mg of dried venom is enough to kill 500 mice subcutaneously and 1,000 intravenously.
P. nigriventer is the species responsible for most cases of envenomation in Brazil because it is commonly found in highly populated areas of southeastern Brazil, such as the states of São Paulo, Minas Gerais, Rio de Janeiro and Espírito Santo.
Reported cases
A 45-year-old man, with no nervous background, employed in the agricultural section of the Butantan Institute, working barefoot, was bitten at 10:40 AM on the small toe of his left foot by a medium-sized Phoneutria. Immediately he felt intense pain that radiated to his foot and leg. He also reported visual disturbances, and when he tried to enter the building, he fell. Without strength, he was supported by two men to the laboratory, where he was examined. He could not stand, had difficulty seeing, difficulty talking (he could not answer questions), and complained of severe general pain and intense cold. He was sweating profusely, there was hyper nasal secretion and salivation, which made him blow his nose and spit constantly. He was agitated, with generalized tremors and continuous cramps in his left foot and leg, and an irregular pulse with 112 beats per minute. An hour after the accident, the serum was injected into the left buttock. During the next hour, the pulse became faster, thready, almost uncountable, and the temperature decreased, with a worsening of the general condition. At noon the victim received a new injection of serum and by 1:00 pm there was improvement.
In another case, a 22-year-old woman, four months pregnant, was bitten on the hand. She experienced severe generalized pain with paroxysm, had difficulty keeping up and talking, moaned, had cramps, generalized tremors, excessive sweating, and rapid pulse (140–150 beats per minute); after receiving the antidote the victim stabilized.
In another case, a 16-year-old boy was bitten on his left hand. The victim had severe general pain, visual disturbance, generalized tremors, cramps, profuse sweating, and a weak, irregular to rapid pulse.
A 23-year-old market worker moving a bunch of bananas was bitten on his hand by P. nigriventer in São Paulo, Brazil. The specimen measured 3.5 cm long and 6 cm with its legs. It was reported that the wound was extremely painful, with the victim noticing that the bite area was sweating and the hair on his skin stood on end. He also reported that the pain radiated to his chest, and that his heart began to race. The victim was dizzy and nauseated, felt cold, began to drool and vomit, and exhibited priapism. He was later treated with anesthetics, tetanus prophylaxis and anti-venom, and recovered 36 hours after the bite.
Another case occurred with a 52-year-old man, bitten by an adult female P. nigriventer. Immediately after the bite he experienced severe local pain, blurred vision, profuse sweating and vomiting. From one to two hours after the bite he presented agitation and high blood pressure; at four hours after the heart rate was high at 150 beats per minute, mild tachypnea, cold extremities, profuse sweating, generalized tremors and priapism. He was treated with anesthetics, anti-venom and fluid replacement.
In 2005, an English man was bitten twice by a spider identified as a Phoneutria, which was hidden in a box of bananas. It was reported that his hand became swollen, he felt dizzy, and that when he got home he collapsed. He was taken to the hospital and received treatment, but his condition continued to deteriorate. He reported chest tightness, difficulty breathing, and both his blood pressure and heart rate were high. He was treated with increased saline to release toxins from the body, and was discharged the next day. He took almost a week to recover.
A 70-year-old man was bitten by a spider with a legspan of 5 to 6 cm. He was cutting sugarcane at the residence of São Pedro de Alcantara, in the countryside, when a spider jumped on his shoulder and came "walking" by the right arm to the back where it bit his hand (bled at the time); he felt intense pain at the time. The spider had several young in the abdomen. The patient evolved with agitation, anxiety, blood pressure 200x110 mmHg, heart rate with 62 beats per minute, respiratory rate 36 breathing movements per minute, respiratory distress, hyperemia, edema and radiating pain. After 50 minutes, the patient was admitted to the emergency department of the HU with blood pressure 150x90 (after captopril), mild dyspnea (patient pneumectomized by TU), with paresthesias and local condition as previously reported. He received anesthetic infiltration, dipyrone and two vials of serum, the right hand remained red and swollen, with improvement in anxiety symptoms and controlled blood pressure.
Another case occurred in Minas Gerais, Brazil, where a man bitten by a Phoneutria developed numbness in his legs, redness, headache and loss of sense of time and space.
Reports of deaths
In a case that occurred in Itanhaém, São Paulo, a 40-year-old man bitten in the foot, presented significant pain and generalized contractures, dying six hours after the accident.
A 7-year-old child, bitten in the ear, presented convulsions, opisthotonos and progressive paralysis, dying 17 hours after the accident.
In a case that occurred in Franca, a 10-year-old child bitten on the middle finger of the right hand, presented severe pain, trismus, tremors in the right arm and face, evolving to permanent contracture, respiratory paralysis, cyanosis and convulsions, dying in 30–40 minutes after the accident.
In São Sebastião, São Paulo, two brothers, 6 months and 18 months old, the children woke up during the night crying and screaming, dying soon after (time of death is not described). The father removed the sheets and found the spider, which was referred to the Butantan Institute, identified as a large female Phoneutria nigriventer.
A 3-year-old girl, bitten on the third finger of her right hand, presented immediate local pain, periods of alternating prostration, cold sweating, chest and abdomen pains and 3 episodes of vomiting, was admitted to UNICAMP 3 hours after the accident, with the same symptoms. 5 vials of AV (Antivenom) and local anesthetic infiltration were administered, and there was improvement in symptoms, with decrease in sweating, but still agitated, 2–3 hours after AV, there was a picture of significant diarrhea (semi-liquid stools) evolving to 2nd degree dehydration. Parenteral hydration was started, 3h30 min post AV, her heart rate was 160 beats per minute, and respiratory rate was 72 at 4 hours post AV. Between 4h15min and 4h30min after AV, she had peripheral cyanosis, a heart rate of 150 beats per minute and dyspnea, and loss of peripheral venous access. Afterwards, there was disseminated pulmonary stertoration, and worsening of agitation, dyspnea and bradycardia, while being performed orotracheal intubation, manual ventilation, adrenaline and external cardiac massage, she later died.
| Biology and health sciences | Spiders | Animals |
277980 | https://en.wikipedia.org/wiki/Man%20o%27%20War | Man o' War | Man o' War (March 29, 1917 – November 1, 1947) was an American Thoroughbred racehorse who is widely regarded as one of the greatest racehorses of all time. Several sports publications, including The Blood-Horse, Sports Illustrated, and the Associated Press, voted Man o' War as the best American racehorse of the 20th century. During his racing career, just after World War I, Man o' War won 20 of 21 races and $249,465 () in purses. He was the unofficial 1920 American horse of the year and was honored with Babe Ruth as the outstanding athlete of the year by The New York Times. He was inducted into the National Museum of Racing and Hall of Fame in 1957. On March 29, 2017, the museum opened a special exhibit in his honor, "Man o' War at 100".
In 1919, Man o' War won 9 of 10 starts, including the Hopeful Stakes and Belmont Futurity, then the most important races for two-year-old horses in the United States. His only loss came at Saratoga Race Course, later nicknamed the Graveyard of Champions, where he lost by a neck to a colt fittingly named Upset.
Man o' War was not entered in the 1920 Kentucky Derby because his owner, Samuel Riddle, did not believe in racing at the distance of so early in a young horse's career. Instead, Man o' War made his three-year-old debut in the Preakness Stakes where he defeated Upset by lengths. Man o' War later won the Belmont Stakes by 20 lengths while setting a world record. Throughout the summer and fall, he continued to dominate his fellow three-year-olds, setting multiple records while conceding large amounts of weight to his rivals. The only time he faced older horses was in the final race of his career in a match race against Sir Barton, who had won what would later be known as the American Triple Crown in 1919. Man o' War won easily by seven lengths in the first horserace to be filmed in its entirety.
Riddle originally intended to race Man o' War in 1921 but decided against it because Man o' War would have been assigned record weights in the handicap format used in almost all races for older horses at the time. Instead, Man o' War was retired to stud, where he became a leading sire whose multiple champions included Triple Crown winner War Admiral. He was the grandsire of Seabiscuit and his sire line continues today through horses such as In Reality, Tiznow, Da' Tara and Tourist. Also successful as a broodmare sire, Man o' War is found in almost all modern American pedigrees.
Background
Man o' War was a chestnut horse with a white star and stripe on his forehead. He was foaled at Nursery Stud near Lexington, Kentucky. He was bred by August Belmont Jr., whose father's accomplishments were recognized through the naming of the Belmont Stakes in 1867. Belmont Jr. was an equally notable horseman who served as the chairman of The Jockey Club from 1895 until his death in 1924. Belmont Park was named in the family's honor when it opened in 1905.
Man o' War was sired by Fair Play, a multiple stakes winner who finished second in the 1908 Belmont Stakes to the undefeated Colin. Man o' War was the second foal out of Mahubah, a lightly raced mare by English Triple Crown Champion Rock Sand. Not long after the colt was foaled on March 29, 1917, Belmont Jr. joined the United States Army at age 65 to serve in France during World War I. While he was overseas, his wife named the foal "Man o' War" in honor of her husband. Originally, the Belmonts intended to race Man o' War themselves. However, in the summer of 1918 with the ongoing war effort, they decided to liquidate their racing stable. At the Saratoga yearling sale, Man o' War was sold at a final bid of $5,000 () to Samuel D. Riddle, who brought him to his Glen Riddle Farm near Berlin, Maryland. The underbidder at the auction was Robert L. Gerry, Sr., who is reported to have said to his wife, "Forty-five hundred is enough to spend for any yearling." Two years later in 1920, Riddle declined an offer of $400,000 for the horse.
At maturity, Man o' War stood , with prominent withers and a high croup, but was sometimes faulted for a dipped back that grew more pronounced with age. He had virtually flawless legs and solid bone, traits he passed on to his offspring. He had a slightly Roman nose and notably high head carriage. His nickname was "Big Red", though his coat had tinges of yellow and gold. An energetic, spirited horse, he is often pictured standing very still and gazing off into the distance, described as the "look of eagles". His stride was measured at 28 feet and, to this day, is believed to still be the longest of all time.
Racing career
Man o' War was the odds-on favorite in every start of his career and justified that faith even in his sole defeat. He typically won in front-running fashion and was only closely pushed in two of his starts. He won the Belmont Stakes by twenty lengths and the Lawrence Realization by a hundred lengths. He set record times in both of those races plus many more at distances ranging from . In many of his starts, he won under heavy restraint and often conceded his rivals large amounts of weight. He retired as the then-leading money-earner in American history.
Developing this talent was not easy for trainer Louis Feustel due to Man o' War's occasionally wild temperament. In his early days, Man o' War would routinely dump his exercise riders, once getting free for over 15 minutes after a morning workout. "He fought like a tiger," Riddle later recalled. "He screamed with rage and fought us so hard that it took several days before he could be handled with safety." Feustel brought the colt along slowly and gradually settled him into a regular routine. Man o' War developed a strong bond with his groom Frank Loftus, who taught the horse to fetch and to carry his hat. Oranges were Man o' War's favorite treat.
1919: Two-year-old season
Man o' War made his debut at Belmont Park on June 6, 1919, in a maiden race over . At the time, horses raced clockwise at Belmont Park rather than counter-clockwise like all the other American tracks of the day, so horses had to learn to race in both directions. This practice ended in 1922 after Man o' War retired. He won his first race by six lengths, then three days later was entered in the Keene Memorial Stakes at a distance of on a muddy track at Belmont Park. At this point in his career, Man o' War had not yet learned how to start quickly but soon settled into fourth place. With an eighth of a mile remaining, jockey Johnny Loftus urged him to run and the horse responded by pulling away from his rivals to win by three lengths. Twelve days later, Man o' War followed up with another win in the Youthful Stakes at Jamaica Race Course. Two days after that, he swept to victory in the Hudson Stakes at Aqueduct to record his fourth win in 18 days. Less than two weeks later, he returned to win the Tremont Stakes at Aqueduct on July 5.
He was then shipped upstate for the summer meet at Saratoga Race Course, where the competition would be much fiercer. He made his first appearance there in the United States Hotel Stakes on August 2. Despite getting a bad start and carrying 130 pounds, Man o' War won by two lengths in his first meeting with a well-regarded colt named Upset.
His next start was the Sanford Memorial Stakes on August 13, where he went off as the odds-on favorite with Upset as the third betting choice. This race is notorious in racing history as Man o' War's only defeat, playing a part in developing Saratoga's reputation as the "Graveyard of Champions". In the early 1900s, there were no starting gates. Horses circled around and then lined up behind a piece of webbing known as the barrier and were sent away when it was raised. In the Sanford, most sources state that Man o' War was still circling with his back to the starting line when the barrier was raised (though some accounts say he was turned only slightly sideways). What is undisputed is that Man o' War had a terrible start, which The New York Times attributed to the absence of the regular starting official. The start was delayed for several minutes as other colts repeatedly broke through the barrier and the starter finally released the field when only the horses near the rail were ready. As a result, Man o' War was far behind the other starters. Loftus then put him in a bad position, getting boxed in by other horses, then checked by tiring horses. Despite this, Man o' War came close to winning, losing by about a neck while conceding 15 pounds to Upset.
It is sometimes erroneously thought that Upset's unexpected win popularized a new phrase in sports (meaning an underdog beating the favorite) – in fact, the term "upset" had been in use to describe such a situation for decades. The two horses faced off five times during their racing careers, with Man o' War beating Upset four times.
The loss only enhanced Man o' War's reputation. J.L. Dempsey of The Daily Racing Form wrote, "Without attempting to detract from the merits of [Upset's] performance, Man o' War proved himself in the running unquestionably the best. It was Upset's advantage at the start, coupled with 15 pounds weight concession, a perfect ride he received from Knapp and his success in saving ground on the stretch turn that brought his triumph over Man o' War. Had the race been a sixteenth farther the finish would have been reversed."
Man o' War gained his revenge with a victory in the Grand Union Hotel Stakes on August 23, beating Upset by two lengths with Blazes in third. After the race, Loftus stated that Man o' War was the best horse he had ever ridden and that his ride had been responsible for the loss in the Sanford. While carrying 130 pounds, Man o' War tied the stakes record of 1:12 2/5 for that had been set by Garbage while carrying only .
Seven days later, Man o' War entered the Hopeful Stakes, whose purse of $30,000 made it one of the richest prizes in racing at the time. He faced eight rivals, a small field at the time for such a prestigious race. Heavy rain started to fall as the field headed to the starting post, and Man o' War broke through the barrier several times, delaying the race by twelve minutes. Despite this, he won with "ridiculous ease" by six lengths.
Man o' War then returned to Belmont Park for the Futurity Stakes on September 14. The Futurity had a purse of $5,000 added, meaning the prize money was increased by nomination and entry fees. Because the Futurity was one of the preeminent races of the day, the added money was large enough to increase the winner's share of the purse to $26,650. Man o' War briefly dueled for the lead with a well-regarded sprinter named Dominique, then opened up a commanding lead. Turning into the stretch, John P. Grier swept into second with an all-out drive but failed to make up any ground on Man o' War, who had not been urged at any point in the race and won by three lengths while carrying 127 pounds, 10 pounds more than John P. Grier. Journalists, horsemen, and fans agreed that Man o' War had to be considered as one of the greatest American horses of his age, and they compared him favorably to the unbeaten Colin. He completed his two-year-old campaign with nine wins from ten starts and earnings of $82,275. He was named the American Champion Two-Year-Old Colt of 1919. He was rated at 136 pounds by handicapper C.C. Ridley of the Daily Racing Form, 16 pounds ahead of the second-ranked colt, Blazes.
1920: Three-year-old season
In 1920, Johnny Loftus was denied a renewal of his jockey's license by The Jockey Club, a development that was rumored to be related to Man o' War's defeat in the Sanford. He was replaced as the colt's rider by Clarence Kummer. Over the winter, Man o' War had grown to high and filled out to about with a girth.
Riddle decided early on to restrict Man o' War to races within his own age division, in large part because the most valuable races in the country were restricted to three-year-olds. He decided not to enter him in the Kentucky Derby because it was run only a few days before his preferred target, the Preakness Stakes, which was held close to the Riddle farm where Man o' War had spent the winter. Riddle also did not like racing in Kentucky and believed it was too early in the year for a young horse to go a mile and a quarter. Thus, Man o' War did not have a chance to complete what later became known as the U.S. Triple Crown of Thoroughbred Racing, consisting of the Kentucky Derby, Preakness Stakes, and Belmont Stakes. The previous year, Sir Barton had won the three races, which gained in prestige and importance 10 years later when Gallant Fox accomplished the same feat under a great deal of media attention.
Instead, Man o' War made his three-year-old debut on May 18 in the Preakness Stakes, then run at a distance of miles. Despite the long layoff and never having raced beyond , he went off as the 4-5 favorite in a field of nine horses that included his old rival Upset, who had finished second in the Derby. Man o' War broke alertly and took the lead within the first 10 yards, then established a two-length lead down the backstretch while under restraint. As they rounded the final turn, Upset started to close ground so Kummer loosened his grip. Man o' War responded by again pulling away, completing the opening mile in what would have been a new Pimlico track record of 1:38. Eased in the final furlong, he won by lengths over Upset in a final time of 1:51.
The horse was then returned to his home base at Belmont Park, where a crowd of 25,000 turned out to watch him in the Withers Stakes on May 29. He was sent off at "generous" odds of 1–7 against two rivals: Wildair, winner of the Metropolitan Handicap who was at odds of 6–1, and the overmatched David Harum at 30–1. Man o' War again seized the early lead, completing the first quarter-mile in : while tugging at the bit. When Wildair tried to close ground around the turn, Kummer briefly released his hold and Man o' War opened up his lead again. Eased in the final sixteenth of a mile, he won by two lengths while setting an American race record of 1:35 for the mile.
Man o' War's next start was on June 12 in the Belmont Stakes, then run at a distance of miles. The race establish him as one of the all-time greats with an "almost unbelievably brilliant performance" before a crowd of 25,000. At odds of 1-20, he faced one rival, the well-regarded colt Donnaconna. Man o' War led from the start, and Kummer let him run freely at the top of the stretch, allowing Man o' War to draw away by 20 lengths. Although eased in the final furlong, he set a world record of 2:14, beating the previous standard set in England by over two seconds and beating Sir Barton's American record by over three seconds. This time stood as the American record until 1961, when Wise Ship ran the distance in 2:14 flat on a turf course. It stood as the American dirt record until 1991.
Ten days later, Man o' War returned in the Stuyvesant Handicap, which he won against one rival. His odds of 1 to 100 were believed to be the lowest ever offered in an American horse race.
His next start in the Dwyer Stakes on July 10 proved far more demanding. The colt John P. Grier, who had challenged Man o' War in the Belmont Futurity at age two, had matured into the second-best three-year-old in the country. Under the conditions of the Dwyer, Man o' War was assigned 126 pounds while John P. Grier carried only 108. The two colts scared away all rivals, turning the Dwyer into a match race. They raced side by side down the backstretch, with Man o' War on the rail blocking John P. Grier from the view of the spectators. The horses ran the race as a sprint, completing the first three-quarters of a mile in 1:09 2/5 – an American record. As they entered the turn, Man o' War started to open up an advantage but John P. Grier rallied and got back on even terms. They completed the mile together in a time of 1:35 3/5, breaking Man o' War's American record set in the Withers. John P. Grier made another surge, and for a moment the spectators believed that he would win the race. Kummer then hit Man o' War with the whip, and he made a final surge, opening up a lead of two lengths in the final fifty yards. The final time was 1:49 1/5, a world record for miles.
Man o' War was then shipped to Saratoga and was entered in the Miller Stakes on August 7. There was a then-record crowd of 35,000, many of whom gathered in the saddling area where Man o' War was surrounded by twelve Pinkerton guards. His jockey for the race was Earl Sande, replacing an injured Kummer. As expected, Man o' War took the early lead and was unchallenged in a six-length victory. He had been tightly restrained, but even so, his time of 1:56 3/5 for the distance of miles was just seconds off the track record. "I never felt anything like that horse in my life," said Sande after the race. "He strides farther than anything I ever rode and does it so handily that you would not think he was running at all."
His connections had a brief scare when Man o' War exhibited signs of lameness after a workout on August 17, but he quickly recovered. On August 21, he entered the Travers Stakes where he faced his two strongest rivals: Upset for the fifth and last time, and John P. Grier for the third time. A record crowd overflowed the grandstand and stretched down the rail and track apron. As a result, Saratoga management opened the infield, and 5,000 people moved across the track to line the inner rail.
Man o' War's jockey for the race was Andy Schuttinger as Kummer was still recovering from injury and Sande was not available. After the start, John P. Grier tried to match strides with Man o' War but could not keep up. After establishing a clear lead by the first turn, Schuttinger took hold of Man o' War for the rest of the race. Upset made a late run to move into second, but Man o' War won by three lengths. Despite not being extended, he completed the distance of miles in 2:01 4/5, equaling the track record set earlier in the year by Sir Barton. This record stood until 1941. It is possible the time was also a world record, as the existing record time of 2:00 flat credited to Whisk Broom II in 1913 was disputed.
Riddle contemplated entering Man o' War against older horses for the first time in the $5,000 Saratoga Gold Cup but ultimately chose to enter the $15,000 Lawrence Realization Stakes against three-year-olds instead. When his only rival scratched from the race, it nearly became a walkover until Mrs. Riddle's niece, Sarah Jeffords, entered Hoodwink. To compensate for the lack of competition, Riddle announced that Man o' War would be allowed to show his speed. He was reunited with regular jockey Kummer, after being ridden by Sande and Schuttinger while Kummer recovered from an injury. Kummer neither restrained the horse nor urged him at any point in the race. Running as he pleased, Man o' War won by slightly more than a quarter-mile – the official margin was 100 lengths – while setting a world record of 2:40 4/5 for a mile and five-eighths. This broke the world record set in England by seconds and lowered the American record by over four seconds. This world record stood until 1956, when it was broken by Man o' War's great-grandson Swaps.
The performance drew superlatives from the racing community. Turf writer B.K. Beckwith later called it "the most astounding display of arrogant annihilation", adding, "[Man o' War] was like a big red sheet of flame running before a prairie wind". The New York Times commented, "at Belmont Park yesterday [Man o' War] gave what was undoubtedly the greatest exhibition of speed over a considerable distance of ground ever witnessed anywhere."
Man o' War's next start was the Jockey Club Gold Cup on September 12, the first time he entered a race open to older horses. The race organizers tried to lure Sir Barton into entering by promising to increase the purse from $15,000 to $25,000 if he did so, but Sir Barton needed more time to recover from his last race. It was also speculated that Sir Barton's owner was holding out for an even higher purse for a match race with Man o' War. The connections of Exterminator, another future Hall of Famer, were also considering the race but bypassed it because, under the weight-for-age conditions of the Gold Cup, he would have had to concede Man o' War five pounds. As it was, Man o' War faced only one competitor and won under tight restraint by fifteen lengths. Although it was declared a hollow victory by The New York Times, Man o' War still set an American record for miles of 2:28 4/5, breaking the existing mark by 4/5 seconds.
Man o' War next entered the Potomac Handicap at Havre de Grace Racetrack in Maryland on September 18. He was assigned 138 pounds, conceding from 24 to 34 pounds to his rivals, which included Kentucky Derby winner Paul Jones. Man o' War faced an early challenge from Blazes, then turned back a late run by Wildair to win by lengths while breaking the track record by seconds. Although Man o' War was not seriously challenged, the high weight and a poorly maintained racing surface took a toll: He came out of the race with a swollen tendon on his right foreleg.
The final start of Man o' War's career came in Windsor, Ontario, Canada, in the Kenilworth Park Gold Cup on October 12. The event was so highly anticipated that it became the first horse race to be filmed in its entirety, with the resulting footage later shown in movie theaters across the country. The race was originally intended to be a face-off between the three great horses of the time: Man o' War, Sir Barton, and Exterminator. However, the owners of Sir Barton and Man o' War agreed to a distance of miles, which was too short for Exterminator to run his best, and they agreed to a weight-for-age format, under which the older Exterminator would concede weight to Man o' War. Therefore, Exterminator was not entered, and in fact, raced that same day on a different track.
In what was now essentially a match race, the advantage shifted to Man o' War, whose front-running style allowed him to dictate the pace. He was almost flat-footed at the break but quickly gathered speed to draw clear of Sir Barton in the first furlong. Kummer then slowed the pace while maintaining a two-length lead down the backstretch. On the final turn, Man o' War put in a brief spurt and quickly opened the lead to five lengths. Kummer again pulled him in and Man o' War maintained a steady pace to the wire to win by seven lengths in a "ridiculously easy" fashion. Although the time of 2:03 was well off the American record shared by Man o' War and Sir Barton, it still broke the track record by over six seconds. Man o' War's share of the purse made him the highest-earning horse in American history. The gold trophy presented in the winner's circle, designed by Tiffany & Co, was later donated by Mrs. Riddle to Saratoga and is now used as the trophy for the Travers Stakes.
There were no formal awards for American Horse of the Year at the time, but Man o' War was informally acknowledged as such in a retroactive poll of turf writers. In its summary of the sporting year, The New York Times stated, "A superman (Babe Ruth) and a superhorse – these were sport's greatest contributions to the history of the year about to close. Some might rate the superhorse, Man o' War, as the outstanding figure of the two, for he has passed on from the field of competition and has left a story of achievement which may never be surpassed."
Over his two-year career, Man o' War won 20 of 21 races, setting three world records, two American records, and three track records.
Weight carrying
At the time, there were few weight-for-age races in the United States. The majority of races were run under handicap conditions, with better horses being assigned higher weights in an attempt to equalize the chances of the other horses in the field. As one writer put it, "the one certain thing about racing is that weight will eventually stop the greatest of racers and bring them down to equal terms with horses of much less quality." As a two-year-old, Man o' War carried in six races; few racehorses have ever carried that much at any age. As a three-year-old, he carried as much as in races, conceding as much as to other horses. Had he raced at age four, the handicappers would have started him at 140 pounds, increasing the amount each time he won until such time as Man o' War either lost or was injured. For this reason, Riddle decided to retire him at the end of his three-year-old campaign.
Race record
Stud record
After his undefeated season as a three-year-old, Man o' War was retired to stud in Lexington, originally at Elizabeth Daingerfield's Hinata Farm. In May 1922, he was moved to Faraway Farm, a property on Huffman Mill Pike outside of Lexington that had been jointly purchased by Riddle, his wife's niece Sarah Jeffords and her husband, Walter. This land, including the stallion barn, is now part of Mt. Brilliant Farm. In late 1936 or early 1937, the Riddle horses were moved to an adjacent property, also called Faraway Farm, where Man o' War spent his final years. This portion of Faraway is now called Man o' War Farm.
During his lifetime, over 500,000 people signed the guest book at Faraway Farm and as many as 1.5 million are estimated to have visited him there. His longtime groom, Will Harbut, would lead the stallion out on display and proudly proclaim his charge's many accomplishments. Man o' War and Harbut were featured together on the cover of the Saturday Evening Post in 1941. This photograph inspired a popular collectors' plate, "Forever Friends", by equine artist Fred Stone.
Man o' War was the leading sire in North America of 1926 and was runner-up in 1928, 1929 and 1937. The Jockey Club credits him with 62 stakes winners from 381 named foals. Riddle restricted Man o' War's book to about 25 mares a year, most of which were owned by him or his friends and family. Though many breeding experts believe that Riddle did not breed the stallion to enough good mares (especially after the first five seasons), he still sired many leading horses. Man o' War's first crop included three horses who were considered champions: American Flag (Belmont Stakes), Florence Nightingale (Coaching Club American Oaks) and Maid at War (Alabama Stakes). His second crop was highlighted by Crusader, who won the Belmont Stakes in 1926 and was largely accepted as the best racehorse of 1926. Among Man o' War's other famous offspring were 1929 Kentucky Derby winner Clyde Van Dusen, Battleship (who won the 1938 English Grand National steeplechase), and War Admiral, the 1937 Triple Crown winner and the second official Horse of the Year. Another of his offspring, Hard Tack, sired Seabiscuit, who was Horse of the Year in 1938.
Man o' War's most successful sons at stud were War Admiral and War Relic, and War Relic's branch of the male line survives today. Tiznow, Tourist, Da' Tara, In Reality, Desert Vixen, Honour and Glory, Bal a Bali, Skywalker and Bertrando are all sire-line descendants of Man o' War. His line is also still active in Europe through classic winner Known Fact and his champion son Warning.
Man o' War's influence in modern pedigrees is magnified by the success of his daughters as broodmares. Although Man o' War never led the broodmare sire list, he was the runner-up ten times and finished in the top ten 22 times. Female line descendants from Man o' War include Eight Thirty, Stymie, Nijinsky, Sword Dancer, Pavot, Riverman, Jim French, Sir Ivor and Kelso. War Admiral was also a leading broodmare sire, especially when crossed with the influential mare La Troienne and his name can be found in many modern pedigrees through such horses as Seattle Slew, Buckpasser and Dr. Fager. American Flag also contributed to Man o' War's modern influence as he was the sire of the second dam of Raise A Native, who is an almost "omnipresent name in American pedigrees". Some modern breeders design matings to concentrate the influence of Man o' War through deep inbreeding. For example, he appears 23 times in the bloodline of American Pharoah, the 2015 Triple Crown winner.
Death
Man o' War was retired from stud in 1943 after suffering a heart attack. He died on November 1, 1947, at age 30 of another apparent heart attack, a short time after Harbut died. His funeral was broadcast live on NBC Radio. Kentucky horseman Ira Drymon said, "He touched the imagination of men and they saw different things in him. But one thing they will all remember was that he brought exaltation into their hearts." An editorial in The New York Times stated: "No other horse ever won such fame as Man o' War. None was more beautiful, with lovelier lines of grace and power. None was more beloved by an admiring and faithful public. Few have lived so long. The American scene seems a little vacant with Man o' War gone to the Elysian Fields where all good horses go."
Man o' War was embalmed and buried in a casket lined with Riddle's black-and-gold racing silks. He was originally interred at Faraway Farm, but in the early 1970s, he was re-interred at a new burial site at the Kentucky Horse Park, where his grave is marked with a statue by American sculptor Herbert Haseltine.
Legacy
Man o' War was inducted into the National Museum of Racing and Hall of Fame in 1957. In 1959, the Man o' War Stakes was created in his honor. In The Blood-Horse magazine 1999 ranking of the top 100 U.S. Thoroughbred champions of the 20th Century, Man o' War was ranked No. 1. He was also ranked No. 1 by the Associated Press as the greatest horse of the 20th century in a separate poll. He was also ranked No. 1 greatest horse in racing history by Sports Illustrated (panel of 7) in 1992. ESPN ranked Man o' War as number 84 on their list of the top North American athletes of the 20th century, compiled for the SportsCentury documentary series. He was one of three racehorses on their list, along with Secretariat (34th) and Citation (97th).
There is a roadside historical marker at the location of his birthplace in Lexington (Nursery Stud), though the site has since been redeveloped. The stallion barn at Faraway Farm where he spent most of his life was renovated in the early 2000s - it is now part of Mt. Brilliant Farm. His stall door was loaned to the National Museum of Racing and Hall of Fame in Saratoga as part a special exhibit, "Man o' War at 100", that opened on March 29, 2017.
Glen Riddle Farm in Maryland was home to Man o' War during the off seasons of his racing career. The land was redeveloped in the early 2000s as GlenRiddle, a gated community with two golf courses, one of which is named in Man o' War's honor. The barn in which Man o' War stayed was converted into the golf clubhouse. In the 1950s, the Riddle estate sold the property Riddle had owned in Glen Riddle, Pennsylvania, which became the Riddlewood housing development. One of the roads in Riddlewood is named Man o' War Drive, another is War Admiral Lane. Man o' War Boulevard, a major arterial that circles to the south of Lexington, is also named in the horse's honor. Central of Georgia Railway once ran a Man o' War passenger train between Atlanta and Columbus.
In books and film
He has been the subject of several biographies. The first, titled Big Red, was by C. W. Anderson and was published in 1943. Man o' War, by Page Cooper and Roger Treat, was published in 1950, and is a classic of its kind. Walter Farley, author of The Black Stallion series, also wrote a slightly fictional biography of Man o' War in 1962. In 2000, Edward L. Bowen wrote a biography called Man o' War: Thoroughbred Legends from Eclipse Press. In 2006, Dorothy Ours wrote a new, extensively sourced biography titled Man o' War: A Legend Like Lightning. Man o' War is also featured in several books on Thoroughbred breeding, including Avalyn Hunter's American Classic Pedigrees (1914–2002).
In 1925, Man o' War was seen in the film Kentucky Pride which was directed by John Ford.
Cultural references
In Sterling A. Brown's poem about Kentucky and an African American in pre-Civil Rights America, "Kentucky Blues" from Southern Road, a reference is made to Man o' War. The poem discusses Thoroughbred horses and other characteristics attributed to the state.
According to Joey DeMaio, the heavy metal band Manowar was named after the horse.
In the 1970 motion picture M*A*S*H, the character Trapper John makes the apocryphal claim that a necropsy performed upon Man o' War following his death revealed that he was a homosexual horse, despite his prolific breeding history.
The song "Action Man" by the band Widespread Panic tells the story of Man o' War. The song references his Sire and Dam, Fair Play and Mahubah, as well as his offspring War Admiral.
Sire line tree
Man o' War
Annapolis
Mercator
American Flag
Gusto
Gun Boat
Crusader
Crossbow II
Royal Crusader
Mars
Scapa Flow
Clyde Van Dusen
Hard Tack
Seabiscuit
Sea Sovereign
Sea Swallow
Battleship
Floating Isle
Navigate
War Battle
Tide Rips
Navy Gun
Sea Legs
Eolus
Cap-A-Pie
Mighty Mo
Shipboard
Fleet Flag
Blockade
Ship Executive
Tsukitomo
War Admiral
Blue Peter
Mr Busher
Cold Command
Navy Page
Admiral Vee
Great War
War Relic
Relic
Buisson Ardent
Roan Rocket
El Relicario
UN Prince
Mystic II
Babamist
Bewley's Hill
Freeman's Hill
Olden Times
Roving Boy
Pericles
Pieces of Eight
Battlefield
Intent
Intentionally
In Reality
Tentam
Group Plan
Missile
Pedigree
Man o' War's grandsire Hastings was purchased by August Belmont Jr. as a two-year-old and won the Belmont Stakes for him in 1896. Hastings had a notoriously vicious temperament that he tended to pass on to his offspring in varying degrees. The mare Fairy Gold, a stakes-winning daughter of Epsom Derby winner Bend Or, was imported by Belmont to America and bred to Hastings in 1904. The resulting foal was Fair Play, who was one of the best racehorses of his year and even better as a stallion, leading the American sire list three times (1920, 1924 and 1927). Fair Play inherited a difficult temperament from Hastings, and his disposition was completely soured after an unsuccessful period spent racing in England. In addition to Man o' War, he also sired Hall of Famer Discovery.
Belmont Jr. also bred the mare Mahubah, who raced only five times with one win. She was the daughter of Rock Sand, who was imported by Belmont to the United States after winning the English Triple Crown. Mahubah produced five foals, all of them by Fair Play, including Man o' War and stakes winner My Play. Mahubah is also the fourth dam of American Triple Crown winner Assault.
An asterisk before a horse's name means the horse was imported to the United States.
Through his sire, Man o' War is a descendant of the first English Triple Crown champion, West Australian. This male line traces to the Godolphin Arabian (Man o' War is 14 generations removed from the Godolphin Arabian).
| Biology and health sciences | Individual animals | Animals |
278003 | https://en.wikipedia.org/wiki/Mountain%20zebra | Mountain zebra | The mountain zebra (Equus zebra) is a zebra species in the family Equidae, native to southwestern Africa. There are two subspecies, the Cape mountain zebra (E. z. zebra) found in South Africa and Hartmann's mountain zebra (E. z. hartmannae) found in south-western Angola and Namibia.
Taxonomy
The mountain zebra comprises two subspecies:
In 2004, C. P. Groves and C. H. Bell investigated the taxonomy of the zebras (genus Equus, subgenus Hippotigris). They concluded that the mountain zebra (Equus zebra zebra) and Hartmann's mountain zebra (Equus zebra hartmannae) are distinct, and suggested that the two would be better classified as separate species, Equus zebra and Equus hartmannae.
However, in a sexual genetic study that included 295 mountain zebra specimens, Moodley and Harley (2005) found nothing to support the separation of the two mountain zebra populations into separate species. They concluded that the Cape mountain zebra and Hartmann's mountain zebra should remain as subspecies.
This is consistent with the third edition of Mammal Species of the World (2005), which lists the mountain zebra as a single species (Equus zebra) with two subspecies.
Appearance
The mountain zebra has a dewlap, which is more conspicuous in E. z. zebra than in E. z. hartmannae. Like all extant zebras, mountain zebras are boldly striped in black or dark brown, and no two individuals look exactly alike. The whole body is striped except for the belly. In the Cape mountain zebra, the ground colour is effectively white, but the ground colour in Hartmann's zebra is slightly buff.
Adult mountain zebras have a head-and-body length of and a tail of long. Wither height ranges from . They weigh from .
Groves and Bell found that Cape mountain zebras exhibit sexual dimorphism, females being larger than males, whereas Hartmann's mountain zebras do not. Hartmann's zebra is on average slightly larger than the Cape mountain zebra.
Habitat
Mountain zebras are found on mountain slopes, open grasslands, woodlands, and areas with sufficient vegetation, but their preferred habitat is mountainous terrain, especially escarpment with a diversity of grass species.
Ecology
Mountain zebras live in hot, dry, rocky, mountainous and hilly habitats. They prefer slopes and plateaus as high as above sea level, although they do migrate lower during winter. Their preferred diet is tufted grass, but in times of shortage, they browse, eating bark, twigs, leaves, buds, fruit, and roots.
They drink every day. When no surface water is available due to drought, they commonly dig for ground water in dry river beds.
The Cape mountain zebra and Hartmann's mountain zebra are now allopatric, meaning that their present ranges do not overlap, which prevents them from crossbreeding. This was not always so, and the current situation is a result of their populations being fragmented when hunters exterminated them throughout the Northern Cape Province of South Africa. Historically, mountain zebras could be found across the entire length of the escarpments along the west coast of southern Africa and in the fold mountain region in the south. However, they generally inhabited poorly productive land and were nowhere really numerous in comparison to those species of zebras or antelope that inhabited the plains, for example.
Behavior and life cycle
Mountain zebras do not aggregate into large herds like plains zebras; they form small family groups consisting of a single stallion and one to five mares, together with their recent offspring. Bachelor males live in separate groups, and mature bachelors attempt to capture young mares to establish a harem. In this they are opposed by the dominant stallion of the group.
Mares give birth to one foal at a time, for about 3 years baby foals gets weaned onto solid forage. Cape mountain zebra foals generally move away from their maternal herds sometime between the ages of 13 and 37 months. However, with Hartmann's mountain zebra, mares try to expel their foals when they are aged around 14 to 16 months. Young males may wander alone for a while before joining a bachelor group, while females are either taken into another breeding herd or are joined by a bachelor male to form a new breeding herd.
Threats
The main threats to the species are the loss of habitat to agriculture, hunting, and persecution. Poaching for food (for example, during guerrilla fighting) has decreased their numbers.
Conservation
The species is listed as vulnerable under the IUCN Red List and both sub-species are included in Appendix II of the Convention on International Trade in Endangered Species (CITES). The Cape mountain zebra was hunted to near extinction. In the 1930s, their population was reduced to about 100 individuals. However, consistent and vigorous conservation measures have succeeded in reversing the decline, and in 1998, the population of the Cape mountain zebra was estimated to have increased to some 1200, about 540 in national parks, 490 in provincial nature reserves, and 165 in other reserves. However, the population has increased to about over 2,700 in the wild due to conservation efforts.
Though both mountain zebra subspecies are currently protected in national parks, they are still threatened. The European Zoos Endangered Species Program and co-operative management of zoo populations worldwide have been set up for them.
| Biology and health sciences | Equidae | Animals |
278366 | https://en.wikipedia.org/wiki/Moment%20%28physics%29 | Moment (physics) | A moment is a mathematical expression involving the product of a distance and a physical quantity such as a force or electric charge. Moments are usually defined with respect to a fixed reference point and refer to physical quantities located some distance from the reference point. For example, the moment of force, often called torque, is the product of a force on an object and the distance from the reference point to the object. In principle, any physical quantity can be multiplied by a distance to produce a moment. Commonly used quantities include forces, masses, and electric charge distributions; a list of examples is provided later.
Elaboration
In its most basic form, a moment is the product of the distance to a point, raised to a power, and a physical quantity
(such as force or electrical charge) at that point:
where is the physical quantity such as a force applied at a point, or a point charge, or a point mass, etc. If the quantity is not concentrated solely at a single point, the moment is the integral of that quantity's density over space:
where is the distribution of the density of charge, mass, or whatever quantity is being considered.
More complex forms take into account the angular relationships between the distance and the physical quantity, but the above equations capture the essential feature of a moment, namely the existence of an underlying or equivalent term. This implies that there are multiple moments (one for each value of n) and that the moment generally depends on the reference point from which the distance is measured, although for certain moments (technically, the lowest non-zero moment) this dependence vanishes and the moment becomes independent of the reference point.
Each value of n corresponds to a different moment: the 1st moment corresponds to n = 1; the 2nd moment to n = 2, etc. The 0th moment (n = 0) is sometimes called the monopole moment; the 1st moment (n = 1) is sometimes called the dipole moment, and the 2nd moment (n = 2) is sometimes called the quadrupole moment, especially in the context of electric charge distributions.
Examples
The moment of force, or torque, is a first moment: , or, more generally, .
Similarly, angular momentum is the 1st moment of momentum: . Momentum itself is not a moment.
The electric dipole moment is also a 1st moment: for two opposite point charges or for a distributed charge with charge density .
Moments of mass:
The total mass is the zeroth moment of mass.
The center of mass is the 1st moment of mass normalized by total mass: for a collection of point masses, or for an object with mass distribution .
The moment of inertia is the 2nd moment of mass: for a point mass, for a collection of point masses, or for an object with mass distribution . The center of mass is often (but not always) taken as the reference point.
Multipole moments
Assuming a density function that is finite and localized to a particular region, outside that region a 1/r potential may be expressed as a series of spherical harmonics:
The coefficients are known as multipole moments, and take the form:
where expressed in spherical coordinates is a variable of integration. A more complete treatment may be found in pages describing multipole expansion or spherical multipole moments. (The convention in the above equations was taken from Jackson – the conventions used in the referenced pages may be slightly different.)
When represents an electric charge density, the are, in a sense, projections of the moments of electric charge: is the monopole moment; the are projections of the dipole moment, the are projections of the quadrupole moment, etc.
Applications of multipole moments
The multipole expansion applies to 1/r scalar potentials, examples of which include the electric potential and the gravitational potential. For these potentials, the expression can be used to approximate the strength of a field produced by a localized distribution of charges (or mass) by calculating the first few moments. For sufficiently large r, a reasonable approximation can be obtained from just the monopole and dipole moments. Higher fidelity can be achieved by calculating higher order moments. Extensions of the technique can be used to calculate interaction energies and intermolecular forces.
The technique can also be used to determine the properties of an unknown distribution . Measurements pertaining to multipole moments may be taken and used to infer properties of the underlying distribution. This technique applies to small objects such as molecules,
but has also been applied to the universe itself, being for example the technique employed by the WMAP and Planck experiments to analyze the cosmic microwave background radiation.
History
In works believed to stem from Ancient Greece, the concept of a moment is alluded to by the word ῥοπή (rhopḗ, "inclination") and composites like ἰσόρροπα (isorropa, "of equal inclinations"). The context of these works is mechanics and geometry involving the lever. In particular, in extant works attributed to Archimedes, the moment is pointed out in phrasings like:
"Commensurable magnitudes ( ) [A and B] are equally balanced () if their distances [to the center Γ, i.e., ΑΓ and ΓΒ] are inversely proportional () to their weights ()."
Moreover, in extant texts such as The Method of Mechanical Theorems, moments are used to infer the center of gravity, area, and volume of geometric figures.
In 1269, William of Moerbeke translates various works of Archimedes and Eutocious into Latin. The term ῥοπή is transliterated into ropen.
Around 1450, Jacobus Cremonensis translates ῥοπή in similar texts into the Latin term momentum ( "movement"). The same term is kept in a 1501 translation by Giorgio Valla, and subsequently by Francesco Maurolico, Federico Commandino, Guidobaldo del Monte, Adriaan van Roomen, Florence Rivault, Francesco Buonamici, Marin Mersenne, and Galileo Galilei. That said, why was the word momentum chosen for the translation? One clue, according to Treccani, is that momento in Medieval Italy, the place the early translators lived, in a transferred sense meant both a "moment of time" and a "moment of weight" (a small amount of weight that turns the scale).
In 1554, Francesco Maurolico clarifies the Latin term momentum in the work Prologi sive sermones. Here is a Latin to English translation as given by Marshall Clagett:
"[...] equal weights at unequal distances do not weigh equally, but unequal weights [at these unequal distances may] weigh equally. For a weight suspended at a greater distance is heavier, as is obvious in a balance. Therefore, there exists a certain third kind of power or third difference of magnitude—one that differs from both body and weight—and this they call moment. Therefore, a body acquires weight from both quantity [i.e., size] and quality [i.e., material], but a weight receives its moment from the distance at which it is suspended. Therefore, when distances are reciprocally proportional to weights, the moments [of the weights] are equal, as Archimedes demonstrated in The Book on Equal Moments. Therefore, weights or [rather] moments like other continuous quantities, are joined at some common terminus, that is, at something common to both of them like the center of weight, or at a point of equilibrium. Now the center of gravity in any weight is that point which, no matter how often or whenever the body is suspended, always inclines perpendicularly toward the universal center.
In addition to body, weight, and moment, there is a certain fourth power, which can be called impetus or force. Aristotle investigates it in On Mechanical Questions, and it is completely different from [the] three aforesaid [powers or magnitudes]. [...]"
in 1586, Simon Stevin uses the Dutch term staltwicht ("parked weight") for momentum in De Beghinselen Der Weeghconst.
In 1632, Galileo Galilei publishes Dialogue Concerning the Two Chief World Systems and uses the Italian momento with many meanings, including the one of his predecessors.
In 1643, Thomas Salusbury translates some of Galilei's works into English. Salusbury translates Latin momentum and Italian momento into the English term moment.
In 1765, the Latin term momentum inertiae (English: moment of inertia) is used by Leonhard Euler to refer to one of Christiaan Huygens's quantities in Horologium Oscillatorium. Huygens 1673 work involving finding the center of oscillation had been stimulated by Marin Mersenne, who suggested it to him in 1646.
In 1811, the French term moment d'une force (English: moment of a force) with respect to a point and plane is used by Siméon Denis Poisson in Traité de mécanique. An English translation appears in 1842.
In 1884, the term torque is suggested by James Thomson in the context of measuring rotational forces of machines (with propellers and rotors). Today, a dynamometer is used to measure the torque of machines.
In 1893, Karl Pearson uses the term n-th moment and in the context of curve-fitting scientific measurements. Pearson wrote in response to John Venn, who, some years earlier, observed a peculiar pattern involving meteorological data and asked for an explanation of its cause. In Pearson's response, this analogy is used: the mechanical "center of gravity" is the mean and the "distance" is the deviation from the mean. This later evolved into moments in mathematics. The analogy between the mechanical concept of a moment and the statistical function involving the sum of the th powers of deviations was noticed by several earlier, including Laplace, Kramp, Gauss, Encke, Czuber, Quetelet, and De Forest.
| Physical sciences | Classical mechanics | Physics |
278490 | https://en.wikipedia.org/wiki/Australian%20funnel-web%20spider | Australian funnel-web spider | Atracidae is a family of mygalomorph spiders, commonly known as Australian funnel-web spiders or atracids. It has been included as a subfamily of the Hexathelidae, but is now recognised as a separate family. All members of the family are native to Australia. Atracidae consists of three genera: Atrax, Hadronyche, and Illawarra, comprising 35 species. Some members of the family produce venom that is dangerous to humans, and bites by spiders of six of the species have caused severe injuries to victims. The bites of the Sydney funnel-web spider (Atrax robustus) and northern tree-dwelling funnel-web spider (Hadronyche formidabilis) are potentially deadly, but no fatalities have occurred since the introduction of modern first-aid techniques and antivenom.
Description
Spiders in the family Atracidae are medium to large in size, with body lengths ranging from 1 to 5 cm (0.4 to 2.0 in), with one exceptional specimen reaching 8 cm (3.1 in). They have a hairless carapace covering the front part of their bodies. Some atracids have relatively long spinnerets; this is especially true of the Sydney funnel-web spider (A. robustus). Males have a large mating spur projecting from the middle of their second pair of legs.
Like other Mygalomorphae – an infraorder of spiders that includes the tropical tarantulas – these spiders have fangs that point straight down the body and do not point towards each other (cf. Araneomorphae). They have ample venom glands that lie entirely within their chelicerae. Their fangs are large and powerful, capable of penetrating fingernails and soft shoes.
Australian funnel-web spiders make their burrows in moist, cool, sheltered habitats – under rocks, in and under rotting logs, and some in rough-barked trees (occasionally meters above ground). They are commonly found in suburban rockeries and shrubberies, rarely in lawns or other open terrain. A burrow characteristically has irregular silk trip-lines radiating from the entrance. Unlike some related trapdoor spiders, they do not build lids to their burrows.
Distribution
The primary range of the Australian funnel-web spiders is the eastern coast of Australia, with specimens found in New South Wales, South Australia, Victoria, Tasmania and Queensland. The only Australian states or territories without members of this family are Western Australia and the Northern Territory.
Taxonomy
The first atracid spider to be described was Hadronyche cerberea, by Carl Ludwig Koch in 1873. Octavius Pickard-Cambridge described another atracid species, A. robustus, four years later. For a considerable time, confusion existed as to the limits of the genera Hadronyche and Atrax, not helped by the destruction of the type specimens of Hadronyche cerberea during World War II. In 1980, Robert J. Raven merged the two genera under Atrax. In 1988, Michael R. Gray separated them again, and in 2010, added a third genus, Illawarra.
The family placement of the group has varied. In 1892, Eugène Simon placed Atrax and Hadronyche in the family Dipluridae. In 1901, Henry R. Hogg considered them to be sufficiently distinctive to form a separate group, which he called "Atraceae" – the basis of the modern family name Atracidae. When in the 1980s, Raven elevated part of Simon's Dipluridae to the family Hexathelidae, he included the atracine group. Molecular phylogenetic studies consistently threw doubt on the monophyly of the Hexathelidae. In 2018, the group was restored to a full family as Atracidae. The following cladogram shows the relationship found between Atracidae and related taxa. Its sister taxon is Actinopodidae.
Genera
, the World Spider Catalog accepts the following genera:
Atrax O. Pickard-Cambridge, 1877 — Australia
Hadronyche L. Koch, 1873 — Australia
Illawarra Gray, 2010 — Australia
Medical significance
Australian funnel-web spiders are one of the most medically significant groups of spiders in the world and are regarded by some to be the most deadly, both in terms of clinical cases and venom toxicity. Six species have caused severe injuries to human victims: the Sydney funnel-web spider (Atrax robustus), the northern tree-dwelling funnel-web spider (Hadronyche formidabilis), the southern tree-dwelling funnel-web spider (H. cerberea), the Blue Mountains funnel-web spider (H. versuta), the Darling Downs funnel-web spider (H. infensa), and the Port Macquarie funnel-web spider (H. macquariensis).
Examination of bite records has implicated wandering males in most, if not all, fatal Australian funnel-web spider bites to humans. Adult males, recognised by the modified terminal segment of the palp, tend to wander during the warmer months of the year, looking for receptive females for mating. They are attracted to water, hence are often found in swimming pools, into which they often fall while wandering. The spiders can survive immersion in water for several hours and can deliver a bite when removed from the water. They also show up in garages and yards in suburban Sydney. Contrary to a commonly held belief, Australian funnel-web spiders are not able to jump, although they can run quickly.
While some very venomous spiders do not always inject venom when they bite, these spiders most often do. The volume of venom delivered to large animals is often small, possibly due to the angle of the fangs, which are not horizontally opposed, and because contact is often brief before the spider is brushed off. About 10 to 25% of bites are claimed to produce significant toxicity, but the likelihood cannot be predicted and all bites should be treated as potentially life-threatening.
Bites from Sydney funnel-web spiders have caused 13 documented deaths (seven in children). In all cases where the sex of the biting spider could be determined, it was found to be the male of the species. One member of the genus Hadronyche, the northern tree-dwelling funnel-web spider, has also been claimed to cause fatal envenomation, but to date, this lacks the support of a specific medical report. Assays of venom from several Hadronyche species have shown it to be similar to Atrax venom.
In 2021 researchers from the University of Queensland announced that the venom from funnel-web spiders found on Fraser Island contains the molecule Hi1a that could block the cell death message after a heart attack, protecting the heart after major injury and potentially preventing or minimising muscle damage.
Toxins
Many different toxins are found in the venom of Atrax and Hadronyche spiders. Collectively, these spider toxins are given the name atracotoxins (ACTX), as all these spiders belong to the family Atracidae. The first toxins isolated were the δ-ACTX toxins present in the venom of both A. robustus (δ-ACTX-Ar1, formerly known as robustoxin or atracotoxin) and H. versuta (δ-ACTX-Hv1a, formerly known as versutoxin). Both of these toxins produce the same effects in monkeys as those seen in humans, suggesting that they are responsible for the physiological effects seen with crude venom.
These toxins are thought to induce spontaneous, repetitive firing and prolongation of action potentials, resulting in continuous acetylcholine neurotransmitter release from somatic and autonomic presynaptic nerve endings. This leads to slower voltage-gated sodium channel inactivation and a hyperpolarizing shift in the voltage-dependence of activation. This in turn inhibits neurally mediated transmitter release, resulting in a surge of endogenous acetylcholine, noradrenaline, and adrenaline.
Although extremely toxic to primates, the venom appears to be fairly harmless to many other animals. These animals may be resistant to the venom's effects due to the presence of IgG, and possibly cross-linked IgG and IgM inactivating factors in their blood plasma that bind to the toxins responsible and neutralise them.
The female's venom was thought to be only about a sixth as potent to humans as that of the male's. The bite of a female or juvenile may still be serious, but considerable variability occurs in venom toxicity between species, together with assumed degrees of inefficiency in the method of venom delivery.
Symptoms
Envenomation symptoms observed following bites by these spiders are very similar. The bite is initially very painful, due to the size of the fangs penetrating the skin. Puncture marks and local bleeding are also usually visible. If substantial envenomation occurs, symptoms generally occur within minutes and progress rapidly.
Early symptoms of systemic envenomation include goose bumps, sweating, tingling around the mouth and tongue, twitching (initially facial and intercostal), salivation, watery eyes, elevated heart rate, and elevated blood pressure. As systemic envenomation progresses, symptoms include nausea, vomiting, shortness of breath (caused by airway obstruction), agitation, confusion, writhing, grimacing, muscle spasms, pulmonary oedema (of neurogenic or hypertensive origin), metabolic acidosis, and extreme hypertension. The final stages of severe envenomation include dilation of the pupils (often fixed), uncontrolled generalised muscle twitching, unconsciousness, elevated intracranial pressure, and death. Death generally is a result of progressive hypotension or possibly elevated intracranial pressure consequent on cerebral oedema.
The onset of severe envenomation can be rapid. In one prospective study, the median time to onset of envenomation was 28 minutes, with only two cases having onset after two hours (both had pressure immobilisation bandages applied). Death may occur within a period ranging from 15 minutes (this occurred when a small child was bitten) to three days.
Treatment
Owing to the severity of symptoms, and the speed with which they progress, in areas where these spiders are known to live, all bites from large, black spiders should be treated as though they were caused by Australian funnel-web spiders. First-aid treatment for a suspected Australian funnel-web spider bite consists of immediately applying a pressure immobilization bandage; a technique which consists of wrapping the bitten limb with a crepe bandage, as well as applying a splint to limit movement of the limb. This technique was originally developed for snakebites, but has also been shown to be effective at slowing venom movement and preventing systemic envenomation in case of an Australian funnel-web spider bite. Some evidence suggests that periods of prolonged localisation may slowly inactivate the venom.
Further supportive care may be necessary, but the mainstay of treatment is antivenom. Venom from the male Sydney funnel-web spider (A. robustus) is used in producing the antivenom, but it appears to be effective against the venom of all species of atracids. Australian funnel-web spider antivenom has also been shown, in vitro, to reverse the effects of eastern mouse spider (Missulena bradleyi) venom.
Before the introduction of antivenom, envenomation resulted in significant morbidity and mortality. The purified rabbit IgG antivenom was developed in 1981 through a team effort led by Dr. Struan Sutherland, head of immunology at the Australian Commonwealth Serum Laboratories in Melbourne. The antivenom is fast-acting and highly and globally effective. Antivenom therapy has shortened the course of envenomation effects; prior to its availability, the average length of hospital treatment for severe bites was about 14 days. Today, antivenom-treated patients are commonly discharged from hospital within one to three days. No deaths are known since it became available.
| Biology and health sciences | Spiders | Animals |
279350 | https://en.wikipedia.org/wiki/Electric%20vehicle | Electric vehicle | An electric vehicle (EV) is a vehicle whose propulsion is powered fully or mostly by electricity. EVs include road and rail vehicles, electric boats and underwater vessels, electric aircraft and electric spacecraft.
Early electric vehicles first came into existence in the late 19th century, when the Second Industrial Revolution brought forth electrification and mass utilization of DC and AC electric motors. Using electricity was among the preferred methods for motor vehicle propulsion as it provides a level of quietness, comfort and ease of operation that could not be achieved by the gasoline engine cars of the time, but range anxiety due to the limited energy storage offered by contemporary battery technologies hindered any mass adoption of private electric vehicles throughout the 20th century. Internal combustion engines (both gasoline and diesel engines) were the dominant propulsion mechanisms for cars and trucks for about 100 years, but electricity-powered locomotion remained commonplace in other vehicle types, such as overhead line-powered mass transit vehicles like electric trains, trams, monorails and trolley buses, as well as various small, low-speed, short-range battery-powered personal vehicles such as mobility scooters.
Hybrid electric vehicles, where electric motors are used as a supplementary propulsion to internal combustion engines, became more widespread in the late 1990s. Plug-in hybrid electric vehicles, where electric motors can be used as the predominant propulsion rather than a supplement, did not see any mass production until the late 2000s, and battery electric cars did not become practical options for the consumer market until the 2010s.
Progress in batteries, electric motors and power electronics have made electric cars more feasible than during the 20th century. As a means of reducing tailpipe emissions of carbon dioxide and other pollutants, and to reduce use of fossil fuels, government incentives are available in many areas to promote the adoption of electric cars and trucks.
History
Electric motive power started in 1827 when Hungarian priest Ányos Jedlik built the first crude but viable electric motor; the next year he used it to power a small model car. In 1835, Professor Sibrandus Stratingh of the University of Groningen, in the Netherlands, built a small-scale electric car, and sometime between 1832 and 1839, Robert Anderson of Scotland invented the first crude electric carriage, powered by non-rechargeable primary cells. American blacksmith and inventor Thomas Davenport built a toy electric locomotive, powered by a primitive electric motor, in 1835. In 1838, a Scotsman named Robert Davidson built an electric locomotive that attained a speed of four miles per hour (6 km/h). In England, a patent was granted in 1840 for the use of rails as conductors of electric current, and similar American patents were issued to Lilley and Colten in 1847.
The first mass-produced electric vehicles appeared in America in the early 1900s. In 1902, the Studebaker Automobile Company entered the automotive business with electric vehicles, though it also entered the gasoline vehicles market in 1904. However, with the advent of cheap assembly line cars by Ford Motor Company, the popularity of electric cars declined significantly.
Due to lack of electricity grids and the limitations of storage batteries at that time, electric cars did not gain much popularity; however, electric trains gained immense popularity due to their economies and achievable speeds. By the 20th century, electric rail transport became commonplace due to advances in the development of electric locomotives. Over time their general-purpose commercial use reduced to specialist roles as platform trucks, forklift trucks, ambulances, tow tractors, and urban delivery vehicles, such as the iconic British milk float. For most of the 20th century, the UK was the world's largest user of electric road vehicles.
Electrified trains were used for coal transport, as the motors did not use the valuable oxygen in the mines. Switzerland's lack of natural fossil resources forced the rapid electrification of their rail network. One of the earliest rechargeable batteriesthe nickel-iron batterywas favored by Edison for use in electric cars.
EVs were among the earliest automobiles, and before the preeminence of light, powerful internal combustion engines (ICEs), electric automobiles held many vehicle land speed and distance records in the early 1900s. They were produced by Baker Electric, Columbia Electric, Detroit Electric, and others, and at one point in history outsold gasoline-powered vehicles. In 1900, 28 percent of the cars on the road in the US were electric. EVs were so popular that even President Woodrow Wilson and his secret service agents toured Washington, D.C., in their Milburn Electrics, which covered 60–70 miles (100–110 km) per charge.
Most producers of passenger cars opted for gasoline cars in the first decade of the 20th century, but electric trucks were an established niche well into the 1920s. Several developments contributed to a decline in the popularity of electric cars. Improved road infrastructure required a greater range than that offered by electric cars, and the discovery of large reserves of petroleum in Texas, Oklahoma, and California led to the wide availability of affordable gasoline/petrol, making internal combustion powered cars cheaper to operate over long distances. Electric vehicles were seldom marketed as women's luxury car, which may have been a stigma among male consumers. Also, internal combustion-powered cars became ever-easier to operate thanks to the invention of the electric starter by Charles Kettering in 1912, which eliminated the need of a hand crank for starting a gasoline engine, and the noise emitted by ICE cars became more bearable thanks to the use of the muffler, which Hiram Percy Maxim had invented in 1897. As roads were improved outside urban areas, the electric vehicle range could not compete with the ICE. Finally, the initiation of mass production of gasoline-powered vehicles by Henry Ford in 1913 reduced significantly the cost of gasoline cars as compared to electric cars.
In the 1930s, National City Lines, which was a partnership of General Motors, Firestone, and Standard Oil of California purchased many electric tram networks across the country to dismantle them and replace them with GM buses. The partnership was convicted of conspiring to monopolize the sale of equipment and supplies to their subsidiary companies. Still, it was acquitted of conspiring to monopolize the provision of transportation services.
The Copenhagen Summit, conducted amid a severe observable climate change brought on by human-made greenhouse gas emissions, was held in 2009. During the summit, more than 70 countries developed plans to reach net zero eventually. For many countries, adopting more EVs will help reduce the use of gasoline.
Experimentation
In January 1990, General Motors President introduced its EV concept two-seater, the "Impact", at the Los Angeles Auto Show. That September, the California Air Resources Board mandated major-automaker sales of EVs, in phases starting in 1998. From 1996 to 1998 GM produced 1117 EV1s, 800 of which were made available through three-year leases.
Chrysler, Ford, GM, Honda, and Toyota also produced limited numbers of EVs for California drivers during this period. In 2003, upon the expiration of GM's EV1 leases, GM discontinued them. The discontinuation has variously been attributed to:
the auto industry's successful federal court challenge to California's zero-emissions vehicle mandate,
a federal regulation requiring GM to produce and maintain spare parts for the few thousand EV1s and
the success of the oil and auto industries' media campaign to reduce public acceptance of EVs.
A movie made on the subject in 2005–2006 was titled Who Killed the Electric Car? and released theatrically by Sony Pictures Classics in 2006. The film explores the roles of automobile manufacturers, oil industry, the U.S. government, batteries, hydrogen vehicles, and the general public, and each of their roles in limiting the deployment and adoption of this technology.
Ford released a number of their Ford Ecostar delivery vans into the market. Honda, Nissan and Toyota also repossessed and crushed most of their EVs, which, like the GM EV1s, had been available only by closed-end lease. After public protests, Toyota sold 200 of its RAV4 EVs; they later sold at over their original forty-thousand-dollar price. Later, BMW of Canada sold off a number of Mini EVs when their Canadian testing ended.
The production of the Citroën Berlingo Electrique stopped in September 2005. Zenn started production in 2006 but ended by 2009.
Reintroduction
During the late 20th and early 21st century, the environmental impact of the petroleum-based transportation infrastructure, along with the fear of peak oil, led to renewed interest in electric transportation infrastructure. EVs differ from fossil fuel-powered vehicles in that the electricity they consume can be generated from a wide range of sources, including fossil fuels, nuclear power, and renewables such as solar power and wind power, or any combination of those. Recent advancements in battery technology and charging infrastructure have addressed many of the earlier barriers to EV adoption, making electric vehicles a more viable option for a wider range of consumers.
The carbon footprint and other emissions of electric vehicles vary depending on the fuel and technology used for electricity generation. The electricity may be stored in the vehicle using a battery, flywheel, or supercapacitors. Vehicles using internal combustion engines usually only derive their energy from a single or a few sources, usually non-renewable fossil fuels. A key advantage of electric vehicles is regenerative braking, which recovers kinetic energy, typically lost during friction braking as heat, as electricity restored to the on-board battery.
Electricity sources
There are many ways to generate electricity, of varying costs, efficiency and ecological desirability.
Connection to generator plants
Direct connection to electric grids as is common among electric trains, trams, trolleybuses, and trolleytrucks ( | Technology | Basics_7 | null |
279624 | https://en.wikipedia.org/wiki/Magnetic%20flux%20quantum | Magnetic flux quantum | The magnetic flux, represented by the symbol , threading some contour or loop is defined as the magnetic field multiplied by the loop area , i.e. . Both and can be arbitrary, meaning that the flux can be as well but increments of flux can be quantized. The wave function can be multivalued as it happens in the Aharonov–Bohm effect or quantized as in superconductors. The unit of quantization is therefore called magnetic flux quantum.
Dirac magnetic flux quantum
The first to realize the importance of the flux quantum was Dirac in his publication on monopoles
The phenomenon of flux quantization was predicted first by Fritz London then within the Aharonov–Bohm effect and later discovered experimentally in superconductors (see below).
Superconducting magnetic flux quantum
If one deals with a superconducting ring (i.e. a closed loop path in a superconductor) or a hole in a bulk superconductor, the magnetic flux threading such a hole/loop is quantized.
The (superconducting) magnetic flux quantum is a combination of fundamental physical constants: the Planck constant and the electron charge .
Its value is, therefore, the same for any superconductor.
To understand this definition in the context of the Dirac flux quantum one shall consider that the effective quasiparticles active in a superconductors are Cooper pairs with an effective charge of 2 electrons .
The phenomenon of flux quantization was first discovered in superconductors experimentally by B. S. Deaver and W. M. Fairbank and, independently, by R. Doll and M. Näbauer, in 1961. The quantization of magnetic flux is closely related to the Little–Parks effect, but was predicted earlier by Fritz London in 1948 using a phenomenological model.
The inverse of the flux quantum, , is called the Josephson constant, and is denoted J. It is the constant of proportionality of the Josephson effect, relating the potential difference across a Josephson junction to the frequency of the irradiation. The Josephson effect is very widely used to provide a standard for high-precision measurements of potential difference, which (from 1990 to 2019) were related to a fixed, conventional value of the Josephson constant, denoted J-90. With the 2019 revision of the SI, the Josephson constant has an exact value of J = .
Derivation of the superconducting flux quantum
The following physical equations use SI units. In CGS units, a factor of would appear.
The superconducting properties in each point of the superconductor are described by the complex quantum mechanical wave function – the superconducting order parameter. As any complex function can be written as , where is the amplitude and is the phase. Changing the phase by will not change and, correspondingly, will not change any physical properties. However, in the superconductor of non-trivial topology, e.g. superconductor with the hole or superconducting loop/cylinder, the phase may continuously change from some value to the value as one goes around the hole/loop and comes to the same starting point. If this is so, then one has magnetic flux quanta trapped in the hole/loop, as shown below:
Per minimal coupling, the current density of Cooper pairs in the superconductor is:
where is the charge of the Cooper pair.
The wave function is the Ginzburg–Landau order parameter:
Plugged into the expression of the current, one obtains:
Inside the body of the superconductor, the current density J is zero, and therefore
Integrating around the hole/loop using Stokes' theorem and gives:
Now, because the order parameter must return to the same value when the integral goes back to the same point, we have:
Due to the Meissner effect, the magnetic induction inside the superconductor is zero. More exactly, magnetic field penetrates into a superconductor over a small distance called London's magnetic field penetration depth (denoted and usually ≈ 100 nm). The screening currents also flow in this -layer near the surface, creating magnetization inside the superconductor, which perfectly compensates the applied field , thus resulting in inside the superconductor.
The magnetic flux frozen in a loop/hole (plus its -layer) will always be quantized. However, the value of the flux quantum is equal to only when the path/trajectory around the hole described above can be chosen so that it lays in the superconducting region without screening currents, i.e. several away from the surface. There are geometries where this condition cannot be satisfied, e.g. a loop made of very thin () superconducting wire or the cylinder with the similar wall thickness. In the latter case, the flux has a quantum different from .
The flux quantization is a key idea behind a SQUID, which is one of the most sensitive magnetometers available.
Flux quantization also plays an important role in the physics of type II superconductors. When such a superconductor (now without any holes) is placed in a magnetic field with the strength between the first critical field and the second critical field , the field partially penetrates into the superconductor in a form of Abrikosov vortices. The Abrikosov vortex consists of a normal core – a cylinder of the normal (non-superconducting) phase with a diameter on the order of the , the superconducting coherence length. The normal core plays a role of a hole in the superconducting phase. The magnetic field lines pass along this normal core through the whole sample. The screening currents circulate in the -vicinity of the core and screen the rest of the superconductor from the magnetic field in the core. In total, each such Abrikosov vortex carries one quantum of magnetic flux .
Measuring the magnetic flux
Prior to the 2019 revision of the SI, the magnetic flux quantum was measured with great precision by exploiting the Josephson effect. When coupled with the measurement of the von Klitzing constant , this provided the most accurate values of the Planck constant obtained until 2019. This may be counterintuitive, since is generally associated with the behaviour of microscopically small systems, whereas the quantization of magnetic flux in a superconductor and the quantum Hall effect are both emergent phenomena associated with thermodynamically large numbers of particles.
As a result of the 2019 revision of the SI, the Planck constant has a fixed value which, together with the definitions of the second and the metre, provides the official definition of the kilogram. Furthermore, the elementary charge also has a fixed value of to define the ampere. Therefore, both the Josephson constant and the von Klitzing constant have fixed values, and the Josephson effect along with the von Klitzing quantum Hall effect becomes the primary mise en pratique for the definition of the ampere and other electric units in the SI.
| Physical sciences | Physical constants | Physics |
280356 | https://en.wikipedia.org/wiki/White-eye | White-eye | The white-eyes are a family, Zosteropidae, of small passerine birds native to tropical, subtropical and temperate Sub-Saharan Africa, southern and eastern Asia, and Australasia. White-eyes inhabit most tropical islands in the Indian Ocean, the western Pacific Ocean, and the Gulf of Guinea. Discounting some widespread members of the genus Zosterops, most species are endemic to single islands or archipelagos. The silvereye, Zosterops lateralis, naturally colonised New Zealand, where it is known as the "wax-eye" or tauhou ("stranger"), from 1855. The silvereye has also been introduced to the Society Islands in French Polynesia, while the Japanese white-eye has been introduced to Hawaii.
Characteristics
White-eyes are mostly of undistinguished appearance, the plumage being generally greenish olive above, and pale grey below. Some species have a white or bright yellow throat, breast, or lower parts, and several have buff flanks. As their common name implies, many species have a conspicuous ring of tiny white feathers around their eyes. The scientific name of the group also reflects this latter feature, being derived from the Ancient Greek for "girdle-eye". They have rounded wings and strong legs. Like many other nectarivorous birds, they have slender, pointed bills, and brush-tipped tongues. The size ranges up to in length.
All the species of white-eyes are sociable, forming large flocks that only separate on the approach of the breeding season. They build trees nests and lay two to four eggs which are usually pale blue. Though mainly insectivorous, they eat nectar and fruits of various kinds. The silvereye can be a problem in Australian vineyards, by piercing the grape allowing infection or insect damage to follow.
Systematics
The family Zosteropidae was introduced (as a subfamily Zosteropinae) in 1853 by the French naturalist Charles Lucien Bonaparte. The white-eyes were long considered a distinct family Zosteropidae because they are rather homogeneous in morphology and ecology, leading to little adaptive radiation and divergence.
The genus Apalopteron, formerly placed in the Meliphagidae, was transferred to the white eyes in 1995 on genetic and behavioral evidence. It differs much in appearance from the typical white-eyes, Zosterops, but is approached by some Micronesian taxa; its color pattern is fairly unusual save the imperfect white eye-ring.
In 2003, Alice Cibois published the results of her study of mtDNA cytochrome b and 12S/16S rRNA sequence data. According to her results, the white-eyes were likely to form a clade also containing the yuhinas, which were until then placed with the Old World babblers, a large "wastebin" family. Previous molecular studies (e.g. Sibley & Ahlquist 1990, Barker et al. 2002) had together with the morphological evidence tentatively placed white-eyes as the Timaliidae's closest relatives already. But some questions remained, mainly because the white-eyes are all very similar birds in habitus and habits, while the Old World babblers are very diverse (because, as we now know, the group as formerly defined was polyphyletic).
Combined with the yuhinas (and possibly other Timaliidae), the limits of the white-eye clade to the "true" Old World babblers becomes indistinct. Therefore, the current (early 2007) opinion weighs towards merging the group into the Timaliidae, perhaps as a subfamily ("Zosteropinae"). Few white-eyes have been thoroughly studied with the new results in mind, however, and almost all of these are from Zosterops which even at this point appears over-lumped. Also, many "Old World babblers" remain in unresolved relationships. Whether there can be a clear delimitation of a white-eye subfamily or even a young or emerging family is a question that requires a more comprehensive study of both this group and Timaliidae to resolve (Jønsson & Fjeldså 2006).
For example, a revision of the yuhinas and the genus Stachyris (Cibois et al. 2002), based on the same genes as Cibois (2003), revealed that the Philippine species placed in the latter genus by some were actually yuhinas. However, when the review by Jønsson & Fjeldså (2006) was published, no study had tried to propose a phylogeny for the newly defined yuhinas including the white eyes. Therefore, Jønsson & Fjeldså (2006) give a rather misleading phylogeny for the group. It appears as if the yuhinas are polyphyletic, with the white-collared yuhina being closer to the ancestor of the Zosterops white-eyes than to other yuhinas including the species moved from Stachyris (Cibois et al. 2002).
In the past, the Madanga (Madanga ruficollis) was included in this family but studies now place it as an atypical member of the Motacillidae.
The cladogram below showing the relationships between families is based on a study of babblers by Tianlong Cai and collaborators published in 2019.
The cladogram below showing the relationships between the genera is based on the study by Carl Oliveros and collaborators that was published in 2021. The genera Apalopteron, Tephrozosterops and Rukia were not sampled in this study. The genus Megazosterops was found to be nested in Heleia. The earlier study by Cai and collaborators found a generally similar phylogeny but with Cleptornis as sister to Heleia. Cai's study found that Apalopteron was nested within Heleia with weak support and that Tephrozosterops was sister to Zosterops.
List of genera
The family contains 149 species divided into 13 genera:
Parayuhina – white-collared yuhina
Staphida – (3 species)
Yuhina – yuhinas (7 species)
Dasycrotapha – babbler and pygmy babblers (3 species)
Sterrhoptilus – babblers (4 species)
Zosterornis – striped babblers (5 species)
Cleptornis – golden white-eye
Rukia – white-eyes (2 species)
Megazosterops – giant white-eye
Heleia – heleias (10 species)
Apalopteron – Bonin white-eye
Tephrozosterops – rufescent darkeye
Zosterops – (110 species including 3 recently extinct)
| Biology and health sciences | Passerida | Animals |
280509 | https://en.wikipedia.org/wiki/Fire%20hydrant | Fire hydrant | A fire hydrant, fireplug, firecock (archaic), hydrant riser or Johnny Pump is a connection point by which firefighters can tap into a water supply. It is a component of active fire protection. Underground fire hydrants have been used in Europe and Asia since at least the 18th century. Above-ground pillar-type hydrants are a 19th-century invention.
Operation
The user (most likely a fire department) attaches a hose to the fire hydrant, then opens a valve on the hydrant to provide a powerful flow of water, on the order of ; this pressure varies according to region and depends on various factors (including the size and location of the attached water main). This user can attach this hose to a fire engine, which can use a powerful pump to boost the water pressure and possibly split it into multiple streams. One may connect the hose with a threaded connection, instantaneous "quick connector" or a Storz connector.
If a fire hydrant is opened or closed too quickly, a water hammer can occur and damage nearby pipes and equipment. The water inside a charged hose line causes it to be very heavy, and high-water pressure causes it to be stiff and unable to make a tight turn while pressurized. When a fire hydrant is unobstructed, this is not a problem, as there is enough room to adequately position the hose.
Most fire hydrant valves are not designed to throttle the water flow; they are designed to be operated full-on or full-off. The valving arrangement of most dry-barrel hydrants is for the drain valve to be open at anything other than full operation. Usage at partial opening can consequently result in considerable flow directly into the soil surrounding the hydrant, which, over time, can cause severe scouring. Gate or butterfly valves can be installed directly onto the hydrant opening to control individual outputs and allow for changing equipment connections without turning off the flow to other outlets. These valves can be up to in diameter to accommodate the large central "steamer" outlets on many US hydrants. It is good practice to install valves on all outlets before using a hydrant as the protective caps are unreliable and can cause major injury if they fail.
New firefighters are often trained extensively on fire hydrants in the fire academy to be quick and safe while connecting the fire engine to the fire hydrant (usually within one minute). Time is often critical as other firefighters will be waiting for the water supply. When operating a hydrant, a firefighter typically wears appropriate personal protective equipment, such as gloves and a helmet with face shield worn. High-pressure water coursing through a potentially aging and corroding hydrant could cause a failure, injuring the firefighter operating the hydrant or bystanders.
In most jurisdictions it is illegal to park a car within a certain distance of a fire hydrant. In North America, the distances are commonly , often indicated by yellow or red paint on the curb. The rationale behind these laws is that hydrants need to be visible and accessible in an emergency. In the event that a car is illegally parked next to a fire hydrant when firefighters need access to it, firefighters are legally allowed to break the car's windows to run the hose through it, while the car owner receives a parking citation.
Other uses
Street pooling
In 1896, during a terrible heatwave in New York City, the Commissioner of Public Works ordered the opening of the fire hydrants to provide relief to the population. Today some US communities provide low flow sprinkler heads to enable residents to use the hydrants to cool off during hot weather, while gaining some control on water usage. Sometimes those simply seeking to play in the water remove the caps and open the valve, providing residents a place to play and cool off in summer.
Preventing misuse
To prevent casual use or misuse, the hydrant requires special tools to be opened, usually a large wrench with a pentagonal socket. Vandals sometimes cause monetary loss by wasting water when they open hydrants. Such vandalism can also reduce municipal water pressure and impair firefighters' efforts to extinguish fires. Most fire hydrants in Australia are protected by a silver-coloured cover with a red top, secured to the ground with bolts to protect the hydrant from vandalism and unauthorized use. The cover must be removed before use.
In most areas of the United States, contractors who need temporary water may purchase permits to use hydrants. The permit will generally require a hydrant meter, a gate valve and sometimes a clapper valve (if not designed into the hydrant already) to prevent backflow into the hydrant. Additionally, residents who wish to use the hydrant to fill their in-ground swimming pool are commonly permitted to do so, provided they pay for the water and agree to allow firefighters to draft from their pool in the case of an emergency.
Municipal services, such as street sweepers and tank trucks, may also be allowed to use hydrants to fill their water tanks. Often sewer maintenance trucks need water to flush out sewerage lines and fill their tanks on site from a hydrant. If necessary, the municipal workers will record the amount of water they used or use a meter.
Fire hydrants may be used to supply water to riot control vehicles. These vehicles use a high-pressure water cannon to discourage rioting.
Since fire hydrants are one of the most accessible parts of a water distribution system, they are often used for attaching pressure gauges or loggers or monitor system water pressure. Automatic flushing devices are often attached to hydrants to maintain chlorination levels in areas of low usage. Hydrants are also used as an easy above-ground access point by leak detection devices to locate leaks from the sound they make.
Construction
Depending on the country or location, hydrants can be above or below ground. In countries including Japan, the UK, Ukraine, Russia or Spain hydrants are accessible under a heavy metal cover. In other countries, such as the US, and many parts of China, an accessible part of the hydrant is above ground. It can also be mounted in an exterior wall of a building.
In areas subject to freezing temperatures, at most only a portion of the hydrant is above ground. The valve is located below the frost line and connected by a riser to the above-ground portion. A valve rod extends from the valve up through a seal at the top of the hydrant, where it can be operated with the proper wrench. This design is known as a "dry barrel" hydrant, in that the barrel, or vertical body of the hydrant, is normally dry. A drain valve underground opens when the water valve is completely closed; this allows all water to drain from the hydrant body to prevent the hydrant from freezing.
In warm areas, above-ground hydrants may be used with one or more valves in the above-ground portion. Unlike with cold-weather hydrants, it is possible to turn the water supply on and off to each port. This style is known as a "wet barrel" hydrant.
Both wet- and dry-barrel hydrants typically have multiple outlets. Wet barrel hydrant outlets are typically individually controlled, while a single stem operates all the outlets of a dry barrel hydrant simultaneously. Thus, wet barrel hydrants allow single outlets to be opened, requiring somewhat more effort, but simultaneously allowing more flexibility.
A typical US dry-barrel hydrant has two smaller outlets and one larger outlet. The larger outlet is often a Storz connection if the local fire department has standardized on hose using Storz fittings for large diameter supply line. The larger outlet is known as a "steamer" connection, because they were once used to supply steam powered water pumps, and a hydrant with such an outlet may be called a "steamer hydrant", although this usage is becoming archaic. Likewise, an older hydrant without a steamer connection may be called a "village hydrant."
Appearance
Above ground hydrants are coloured for purely practical criteria or more aesthetic reasons. In the United States, the AWWA and NFPA recommend hydrants be colored chrome yellow for rapid identification apart from the bonnet and nozzle caps which should be coded according to their available flow. Class AA hydrants (>1500 gpm) should have their nozzle caps and bonnet colored light blue, Class A hydrants (1000–1499 gpm) green, Class B hydrants (500–999 gpm) orange, Class C hydrants (0–499 gpm) red, and inoperable or end-of-system (risking water hammer) black. These aids arriving firefighters in determining how much water is available and whether to call for additional resources or find another hydrant.
Other codings can be and frequently are used, some of greater complexity, incorporating pressure information, others more simplistic. In Ottawa, Ontario, hydrant colors communicate different messages to firefighters; for example, if the inside of the hydrant is corroded so much that the interior diameter is too narrow for good pressure, it will be painted in a specific scheme to indicate to firefighters to move on to the next one. In many localities, a white or purple top indicates that the hydrant provides non-potable water. Where artistic and/or aesthetic considerations are paramount, hydrants can be extremely varied, or more subdued. In both instances this is usually at the cost of reduced practicality.
In Germany, the Netherlands, Spain, the UK, and many other countries, most hydrants are located below ground and are reached by a riser, which provides the connections for the hoses. The covers can also be artistically designed.
Signage
In the United Kingdom and Ireland, hydrants are located in the ground. Yellow "H" hydrant signs indicate the location of the hydrants and are similar to the blue signs in Finland. Mounted on a small post or nearby wall etc., the two numbers indicate the diameter of the water main (top number) and the distance from the sign (lower number). Modern signs show these measurements in millimetres and metres, whereas older signs use inches and feet. Because the orders of magnitude are so different (6 inches versus 150 mm) there is no ambiguity whichever measuring system is used.
In areas of the United States without winter snow cover, blue reflectors embedded in the street are used to allow rapid identification of hydrants at night. In areas with snow cover, tall signs or flags are used so that hydrants can be found even if covered with snow. In rural areas tall narrow posts painted with visible colours such as red are attached to the hydrants to allow them to be found during heavy snowfall periods. The tops of the fire hydrants indicate available flow in gallons per minute; the color helps make a more accurate choice of what hydrants will be utilized to supply water to the fire scene.
Blue: or more; very good flow
Green: ; good for residential areas
Orange: ; marginally adequate
Red: below ; inadequate
The hydrant bodies are also color-coded.
Chrome Yellow: Municipal System
Red: Private System
Violet: Non-potable supply
These markings and colours are prescribed in NFPA 291: Recommended Practice for Water Flow Testing and Marking of Hydrants. but most municipal water authorities do not actually follow these guidelines.
In Australia, hydrant signage varies, with several types displayed across the country. Most Australian hydrants are underground, being of a ballcock system (spring hydrant type), and a separate standpipe with a central plunger is used to open the valve. Consequently, hydrant signage is essential, because of their concealed nature.
Painted markersUsually a white or yellow (sometimes reflective paint) triangle or arrow painted on the road, pointing towards the side of the road the hydrant will be found on. These are most common in old areas, or on new roads where more advanced signs have not been installed. These are almost always coupled with a secondary form of signage.
Hydrant Marker PlatesFound on power poles, fences, or street-signs, these are a comprehensive and effective system of identification. The plate consists of several codes; H (Potable water Hydrant), RH (Recycled/Non-Potable), P (Pathway, where the hydrant cover can be found), R (Roadway). The plate is vertically oriented, around 8 cm wide, and 15 cm high. It usually faces in the direction of the hydrant. Found on this plate, from top to bottom, are the following features:
The codes listed above, Potable/Non-potable at the top, Path/Roadway on the bottom of the plate.
Below this, a number giving the distance to the hydrant (in meters), then a second number below that giving the size (in millimeters) of the water main.
A black line across the center of the plate indicates the hydrant is found on the opposite side of the road to which the plate is affixed.
Plates for recycled water have a purple background, as well as the RH code, normal potable hydrants are white, with the H code.
Road markers or Cat's eyesAlmost exclusively blue, these are placed on one side or the other of the centre line of the road, to indicate on which side of the road the hydrant lies. They are visible for several hundred meters at night in heavy rain, further in clear conditions.
In Germany the hydrant marker plates follow the style of other marker plates pointing to underground installations. Fire hydrant marker plates have a red border. Other water hydrants may have a blue border. A gas hydrant would have a yellow background instead of a white one for fire hydrants. All of them have large central T with the installation identification on top of itan "H" or older "UH" is located in the ground, a "OH" is above ground, followed by the pipe inner diameter in millimeters (with a small 80 mm in residential areas). The numbers around the T allow to locate the installation in reference to the plate's locationthe number left of the T is in meter left of the sign, the number right of the T is in meter right of the sign, and number below the T tells the distance in meter in front of the sign, where a negative number would point to a place behind the sign. The distance numbers are always given with a comma decimeter precision. If it is not a common fire hydrant type then another identification may be used, for example "300 m³" would point to a cistern usable to pump water from.
In East Asia (China, Japan and South Korea) and former Socialist countries of Eastern Europe, there are two types of fire hydrants, of which one is on the public ground and the other inside a building. The ones inside a building are installed on a wall. They are big, rectangular boxes that also provide alarms (sirens), a fire extinguisher and, at certain times, emergency kits.
Inspection and maintenance
In most areas, fire hydrants require annual inspections and maintenance; they normally only have a one-year warranty, but some have 5- or even 10-year warranties, although the longer warranty does not remove the need for periodic inspections or maintenance. These inspections are generally performed by the local municipalities or fire departments, but they often do not inspect hydrants that are identified as private. Private hydrants are usually located on larger properties to adequately protect large buildings in case of a fire and in order to comply with the fire code. Such hydrants have met the requirements of insurance underwriters and are often referred to as UL/FM hydrants. Usually, companies are contracted out to inspect private fire hydrants, unless the municipality has undertaken that task.
Some fire hydrant manufacturers recommend lubricating the head mechanism and restoring the head gaskets and O-rings annually in order that the fire hydrant perform the service expected of them, while others have incorporated proprietary features to provide long-term lubrication of the hydrant's operating mechanism. In any case, periodic inspection of lubricants is recommended. Lubrication is generally done with a food-grade non-petroleum lubricant to avoid contamination of the distribution system.
Occasionally a stone or foreign object will mar the seat gasket. In this case, most hydrants have a special seat wrench that allows removal of the seat to replace the gasket or other broken parts without removing the hydrant from the ground. Hydrant extensions are also available for raising a hydrant if the grade around the hydrant changes. Without extending the height, the wrenches to remove caps would not clear and the break flanges for traffic models would not be located correctly in case they were hit. Hydrant repair kits are also available to repair sacrificial parts designed to break when hit by a vehicle.
Many fire departments use the hydrants for flushing out water line sediments. When doing so, they often use a hydrant diffuser, a device that diffuses the water so that it does not damage property and is less dangerous to bystanders than a solid stream. Some diffusers also dechlorinate the water to avoid ground contamination. Hydrants are also sometimes used as entry or exit points for pipe cleaning pigs.
In 2011, Code for America developed an "Adopt a Hydrant" website, which enables volunteers to sign up to shovel out fire hydrants after snowstorms. As of 2014, the system has been implemented in Boston; Providence, Rhode Island; Anchorage, Alaska; and Chicago.
Non-pressurized (dry) Hydrants
In rural areas where municipal water systems are not available, dry hydrants are used to supply water for fighting fires. A dry hydrant is analogous to a standpipe. A dry hydrant is usually an unpressurized, permanently installed pipe that has one end below the water level of a lake or pond. This end usually has a strainer to prevent debris or wildlife, such as fish, from entering the pipe. The other end is above ground and has a hard sleeve connector.
When needed, a pumper fire engine will pump from the lake or pond by drafting water. This is done by vacuuming the air out of the dry hydrant, hard sleeve, and the fire engine pump with a primer. Because lower pressure now exists at the pump intake, atmospheric pressure on the water and the weight of the water forces water into the above-water portion of the dry hydrant, into the hard sleeve, and finally into the pump. This water can then be pumped by the engine's centrifugal pump.
Other types
Water wells are also sometimes classified as fire hydrants if they can supply enough water volume and pressure.
Standpipes are connections for firehoses within a building and serve the same purpose inside larger structures as fire hydrants do outdoors. Standpipes may be "dry" or "wet" (permanently filled with water); a dry standpipe requires an external source of water such as firefighting equipment.
History
Before piped mains supplies, water for firefighting had to be kept in buckets and cauldrons ready for use by 'bucket-brigades' or brought with a horse-drawn fire-pump. From the 16th century, as wooden mains water systems were installed, firefighters would dig down to the pipes and drill a hole for water to fill a “wet well” for the buckets or pumps. This had to be filled and plugged afterwards, hence the common US term for a hydrant, 'fireplug'. A marker would be left to indicate where a 'plug' had already been drilled to enable firefighters to find ready-drilled holes. Later wooden systems had pre-drilled holes and plugs.
When cast iron pipes replaced the wood, permanent underground access points were included for the fire fighters. Some countries provide access covers to these points, while others attach fixed above-ground hydrantsthe first cast iron ones were patented in 1801 by Frederick Graff, then chief-engineer of the Philadelphia Water Works. Invention since then has targeted problems such as tampering, freezing, connection, reliability etc.
| Technology | Fire protection | null |
280582 | https://en.wikipedia.org/wiki/Numerical%20digit | Numerical digit | A numerical digit (often shortened to just digit) or numeral is a single symbol used alone (such as "1"), or in combinations (such as "15"), to represent numbers in positional notation, such as the common base 10. The name "digit" originates from the Latin digiti meaning fingers.
For any numeral system with an integer base, the number of different digits required is the absolute value of the base. For example, decimal (base 10) requires ten digits (0 to 9), and binary (base 2) requires only two digits (0 and 1). Bases greater than 10 require more than 10 digits, for instance hexadecimal (base 16) requires 16 digits (usually 0 to 9 and A to F).
Overview
In a basic digital system, a numeral is a sequence of digits, which may be of arbitrary length. Each position in the sequence has a place value, and each digit has a value. The value of the numeral is computed by multiplying each digit in the sequence by its place value, and summing the results.
Digital values
Each digit in a number system represents an integer. For example, in decimal the digit "1" represents the integer one, and in the hexadecimal system, the letter "A" represents the number ten. A positional number system has one unique digit for each integer from zero up to, but not including, the radix of the number system.
Thus in the positional decimal system, the numbers 0 to 9 can be expressed using their respective numerals "0" to "9" in the rightmost "units" position. The number 12 is expressed with the numeral "2" in the units position, and with the numeral "1" in the "tens" position, to the left of the "2" while the number 312 is expressed with three numerals: "3" in the "hundreds" position, "1" in the "tens" position, and "2" in the "units" position.
Computation of place values
The decimal numeral system uses a decimal separator, commonly a period in English, or a comma in other European languages, to denote the "ones place" or "units place", which has a place value one. Each successive place to the left of this has a place value equal to the place value of the previous digit times the base. Similarly, each successive place to the right of the separator has a place value equal to the place value of the previous digit divided by the base. For example, in the numeral 10.34 (written in base 10),
the 0 is immediately to the left of the separator, so it is in the ones or units place, and is called the units digit or ones digit;
the 1 to the left of the ones place is in the tens place, and is called the tens digit;
the 3 is to the right of the ones place, so it is in the tenths place, and is called the tenths digit;
the 4 to the right of the tenths place is in the hundredths place, and is called the hundredths digit.
The total value of the number is 1 ten, 0 ones, 3 tenths, and 4 hundredths. The zero, which contributes no value to the number, indicates that the 1 is in the tens place rather than the ones place.
The place value of any given digit in a numeral can be given by a simple calculation, which in itself is a complement to the logic behind numeral systems. The calculation involves the multiplication of the given digit by the base raised by the exponent , where n represents the position of the digit from the separator; the value of n is positive (+), but this is only if the digit is to the left of the separator. And to the right, the digit is multiplied by the base raised by a negative (−) n. For example, in the number 10.34 (written in base 10),
the 1 is second to the left of the separator, so based on calculation, its value is,
the 4 is second to the right of the separator, so based on calculation its value is,
History
The first true written positional numeral system is considered to be the Hindu–Arabic numeral system. This system was established by the 7th century in India, but was not yet in its modern form because the use of the digit zero had not yet been widely accepted. Instead of a zero sometimes the digits were marked with dots to indicate their significance, or a space was used as a placeholder. The first widely acknowledged use of zero was in 876. The original numerals were very similar to the modern ones, even down to the glyphs used to represent digits.
By the 13th century, Western Arabic numerals were accepted in European mathematical circles (Fibonacci used them in his ). They began to enter common use in the 15th century. By the end of the 20th century virtually all non-computerized calculations in the world were done with Arabic numerals, which have replaced native numeral systems in most cultures.
Other historical numeral systems using digits
The exact age of the Maya numerals is unclear, but it is possible that it is older than the Hindu–Arabic system. The system was vigesimal (base 20), so it has twenty digits. The Mayas used a shell symbol to represent zero. Numerals were written vertically, with the ones place at the bottom. The Mayas had no equivalent of the modern decimal separator, so their system could not represent fractions.
The Thai numeral system is identical to the Hindu–Arabic numeral system except for the symbols used to represent digits. The use of these digits is less common in Thailand than it once was, but they are still used alongside Arabic numerals.
The rod numerals, the written forms of counting rods once used by Chinese and Japanese mathematicians, are a decimal positional system able to represent not only zero but also negative numbers. Counting rods themselves predate the Hindu–Arabic numeral system. The Suzhou numerals are variants of rod numerals.
Modern digital systems
In computer science
The binary (base 2), octal (base 8), and hexadecimal (base 16) systems, extensively used in computer science, all follow the conventions of the Hindu–Arabic numeral system. The binary system uses only the digits "0" and "1", while the octal system uses the digits from "0" through "7". The hexadecimal system uses all the digits from the decimal system, plus the letters "A" through "F", which represent the numbers 10 to 15 respectively. When the binary system is used, the term "bit(s)" is typically used as an alternative for "digit(s)", being a portmanteau of the term "binary digit".
Unusual systems
The ternary and balanced ternary systems have sometimes been used. They are both base 3 systems.
Balanced ternary is unusual in having the digit values 1, 0 and –1. Balanced ternary turns out to have some useful properties and the system has been used in the experimental Russian Setun computers.
Several authors in the last 300 years have noted a facility of positional notation that amounts to a modified decimal representation. Some advantages are cited for use of numerical digits that represent negative values. In 1840 Augustin-Louis Cauchy advocated use of signed-digit representation of numbers, and in 1928 Florian Cajori presented his collection of references for negative numerals. The concept of signed-digit representation has also been taken up in computer design.
Digits in mathematics
Despite the essential role of digits in describing numbers, they are relatively unimportant to modern mathematics. Nevertheless, there are a few important mathematical concepts that make use of the representation of a number as a sequence of digits.
Digital roots
The digital root is the single-digit number obtained by summing the digits of a given number, then summing the digits of the result, and so on until a single-digit number is obtained.
Casting out nines
Casting out nines is a procedure for checking arithmetic done by hand. To describe it, let represent the digital root of , as described above. Casting out nines makes use of the fact that if , then . In the process of casting out nines, both sides of the latter equation are computed, and if they are not equal, the original addition must have been faulty.
Repunits and repdigits
Repunits are integers that are represented with only the digit 1. For example, 1111 (one thousand, one hundred and eleven) is a repunit. Repdigits are a generalization of repunits; they are integers represented by repeated instances of the same digit. For example, 333 is a repdigit. The primality of repunits is of interest to mathematicians.
Palindromic numbers and Lychrel numbers
Palindromic numbers are numbers that read the same when their digits are reversed. A Lychrel number is a positive integer that never yields a palindromic number when subjected to the iterative process of being added to itself with digits reversed. The question of whether there are any Lychrel numbers in base 10 is an open problem in recreational mathematics; the smallest candidate is 196.
History of ancient numbers
Counting aids, especially the use of body parts (counting on fingers), were certainly used in prehistoric times as today. There are many variations. Besides counting ten fingers, some cultures have counted knuckles, the space between fingers, and toes as well as fingers. The Oksapmin culture of New Guinea uses a system of 27 upper body locations to represent numbers.
To preserve numerical information, tallies carved in wood, bone, and stone have been used since prehistoric times. Stone age cultures, including ancient indigenous American groups, used tallies for gambling, personal services, and trade-goods.
A method of preserving numeric information in clay was invented by the Sumerians between 8000 and 3500 BC. This was done with small clay tokens of various shapes that were strung like beads on a string. Beginning about 3500 BC, clay tokens were gradually replaced by number signs impressed with a round stylus at different angles in clay tablets (originally containers for tokens) which were then baked. About 3100 BC, written numbers were dissociated from the things being counted and became abstract numerals.
Between 2700 and 2000 BC, in Sumer, the round stylus was gradually replaced by a reed stylus that was used to press wedge-shaped cuneiform signs in clay. These cuneiform number signs resembled the round number signs they replaced and retained the additive sign-value notation of the round number signs. These systems gradually converged on a common sexagesimal number system; this was a place-value system consisting of only two impressed marks, the vertical wedge and the chevron, which could also represent fractions. This sexagesimal number system was fully developed at the beginning of the Old Babylonia period (about 1950 BC) and became standard in Babylonia.
Sexagesimal numerals were a mixed radix system that retained the alternating base 10 and base 6 in a sequence of cuneiform vertical wedges and chevrons. By 1950 BC, this was a positional notation system. Sexagesimal numerals came to be widely used in commerce, but were also used in astronomical and other calculations. This system was exported from Babylonia and used throughout Mesopotamia, and by every Mediterranean nation that used standard Babylonian units of measure and counting, including the Greeks, Romans and Egyptians. Babylonian-style sexagesimal numeration is still used in modern societies to measure time (minutes per hour) and angles (degrees).
History of modern numbers
In China, armies and provisions were counted using modular tallies of prime numbers. Unique numbers of troops and measures of rice appear as unique combinations of these tallies. A great convenience of modular arithmetic is that it is easy to multiply. This makes use of modular arithmetic for provisions especially attractive. Conventional tallies are quite difficult to multiply and divide. In modern times modular arithmetic is sometimes used in digital signal processing.
The oldest Greek system was that of the Attic numerals, but in the 4th century BC they began to use a quasidecimal alphabetic system (see Greek numerals). Jews began using a similar system (Hebrew numerals), with the oldest examples known being coins from around 100 BC.
The Roman empire used tallies written on wax, papyrus and stone, and roughly followed the Greek custom of assigning letters to various numbers. The Roman numerals system remained in common use in Europe until positional notation came into common use in the 16th century.
The Maya of Central America used a mixed base 18 and base 20 system, possibly inherited from the Olmec, including advanced features such as positional notation and a zero. They used this system to make advanced astronomical calculations, including highly accurate calculations of the length of the solar year and the orbit of Venus.
The Incan Empire ran a large command economy using quipu, tallies made by knotting colored fibers. Knowledge of the encodings of the knots and colors was suppressed by the Spanish conquistadors in the 16th century, and has not survived although simple quipu-like recording devices are still used in the Andean region.
Some authorities believe that positional arithmetic began with the wide use of counting rods in China. The earliest written positional records seem to be rod calculus results in China around 400. Zero was first used in India in the 7th century CE by Brahmagupta.
The modern positional Arabic numeral system was developed by mathematicians in India, and passed on to Muslim mathematicians, along with astronomical tables brought to Baghdad by an Indian ambassador around 773.
From India, the thriving trade between Islamic sultans and Africa carried the concept to Cairo. Arabic mathematicians extended the system to include decimal fractions, and Muḥammad ibn Mūsā al-Ḵwārizmī wrote an important work about it in the 9th century. The modern Arabic numerals were introduced to Europe with the translation of this work in the 12th century in Spain and Leonardo of Pisa's Liber Abaci of 1201. In Europe, the complete Indian system with the zero was derived from the Arabs in the 12th century.
The binary system (base 2) was propagated in the 17th century by Gottfried Leibniz. Leibniz had developed the concept early in his career, and had revisited it when he reviewed a copy of the I Ching from China. Binary numbers came into common use in the 20th century because of computer applications.
Numerals in most popular systems
Additional numerals
| Mathematics | Basics | null |
280911 | https://en.wikipedia.org/wiki/Confidence%20interval | Confidence interval | Informally, in frequentist statistics, a confidence interval (CI) is an interval which is expected to typically contain the parameter being estimated. More specifically, given a confidence level (95% and 99% are typical values), a CI is a random interval which contains the parameter being estimated % of the time. The confidence level, degree of confidence or confidence coefficient represents the long-run proportion of CIs (at the given confidence level) that theoretically contain the true value of the parameter; this is tantamount to the nominal coverage probability. For example, out of all intervals computed at the 95% level, 95% of them should contain the parameter's true value.
Factors affecting the width of the CI include the sample size, the variability in the sample, and the confidence level. All else being the same, a larger sample produces a narrower confidence interval, greater variability in the sample produces a wider confidence interval, and a higher confidence level produces a wider confidence interval.
History
Methods for calculating confidence intervals for the binomial proportion appeared from the 1920s. The main ideas of confidence intervals in general were developed in the early 1930s, and the first thorough and general account was given by Jerzy Neyman in 1937.
Neyman described the development of the ideas as follows (reference numbers have been changed):
[My work on confidence intervals] originated about 1930 from a simple question of Waclaw Pytkowski, then my student in Warsaw, engaged in an empirical study in farm economics. The question was: how to characterize non-dogmatically the precision of an estimated regression coefficient? ...
Pytkowski's monograph ... appeared in print in 1932. It so happened that, somewhat earlier, Fisher published his first paper concerned with fiducial distributions and fiducial argument. Quite unexpectedly, while the conceptual framework of fiducial argument is entirely different from that of confidence intervals, the specific solutions of several particular problems coincided. Thus, in the first paper in which I presented the theory of confidence intervals, published in 1934, I recognized Fisher's priority for the idea that interval estimation is possible without any reference to Bayes' theorem and with the solution being independent from probabilities a priori. At the same time I mildly suggested that Fisher's approach to the problem involved a minor misunderstanding.
In medical journals, confidence intervals were promoted in the 1970s but only became widely used in the 1980s. By 1988, medical journals were requiring the reporting of confidence intervals.
Definition
Let be a random sample from a probability distribution with statistical parameter , which is a quantity to be estimated, and , representing quantities that are not of immediate interest. A confidence interval for the parameter , with confidence level or coefficient , is an interval determined by random variables and with the property:
The number , whose typical value is close to but not greater than 1, is sometimes given in the form (or as a percentage ), where is a small positive number, often 0.05.
It is important for the bounds and to be specified in such a way that as long as is collected randomly, every time we compute a confidence interval, there is probability that it would contain , the true value of the parameter being estimated. This should hold true for any actual and .
Approximate confidence intervals
In many applications, confidence intervals that have exactly the required confidence level are hard to construct, but approximate intervals can be computed. The rule for constructing the interval may be accepted as providing a confidence interval at level if
to an acceptable level of approximation. Alternatively, some authors simply require that
which is useful if the probabilities are only partially identified or imprecise, and also when dealing with discrete distributions. Confidence limits of the form
and
are called conservative; accordingly, one speaks of conservative confidence intervals and, in general, regions.
Desired properties
When applying standard statistical procedures, there will often be standard ways of constructing confidence intervals. These will have been devised so as to meet certain desirable properties, which will hold given that the assumptions on which the procedure relies are true. These desirable properties may be described as: validity, optimality, and invariance.
Of the three, "validity" is most important, followed closely by "optimality". "Invariance" may be considered as a property of the method of derivation of a confidence interval, rather than of the rule for constructing the interval. In non-standard applications, these same desirable properties would be sought:
Validity
This means that the nominal coverage probability (confidence level) of the confidence interval should hold, either exactly or to a good approximation.
Optimality
This means that the rule for constructing the confidence interval should make as much use of the information in the data-set as possible.
One way of assessing optimality is by the width of the interval so that a rule for constructing a confidence interval is judged better than another if it leads to intervals whose widths are typically shorter.
Invariance
In many applications, the quantity being estimated might not be tightly defined as such.
For example, a survey might result in an estimate of the median income in a population, but it might equally be considered as providing an estimate of the logarithm of the median income, given that this is a common scale for presenting graphical results. It would be desirable that the method used for constructing a confidence interval for the median income would give equivalent results when applied to constructing a confidence interval for the logarithm of the median income: Specifically the values at the ends of the latter interval would be the logarithms of the values at the ends of former interval.
Methods of derivation
For non-standard applications, there are several routes that might be taken to derive a rule for the construction of confidence intervals. Established rules for standard procedures might be justified or explained via several of these routes. Typically a rule for constructing confidence intervals is closely tied to a particular way of finding a point estimate of the quantity being considered.
Summary statistics
This is closely related to the method of moments for estimation. A simple example arises where the quantity to be estimated is the population mean, in which case a natural estimate is the sample mean. Similarly, the sample variance can be used to estimate the population variance. A confidence interval for the true mean can be constructed centered on the sample mean with a width which is a multiple of the square root of the sample variance.
Likelihood theory
Estimates can be constructed using the maximum likelihood principle, the likelihood theory for this provides two ways of constructing confidence intervals or confidence regions for the estimates.
Estimating equations
The estimation approach here can be considered as both a generalization of the method of moments and a generalization of the maximum likelihood approach. There are corresponding generalizations of the results of maximum likelihood theory that allow confidence intervals to be constructed based on estimates derived from estimating equations.
Hypothesis testing
If hypothesis tests are available for general values of a parameter, then confidence intervals/regions can be constructed by including in the confidence region all those points for which the hypothesis test of the null hypothesis that the true value is the given value is not rejected at a significance level of
Bootstrapping
In situations where the distributional assumptions for the above methods are uncertain or violated, resampling methods allow construction of confidence intervals or prediction intervals. The observed data distribution and the internal correlations are used as the surrogate for the correlations in the wider population.
Central limit theorem
The central limit theorem is a refinement of the law of large numbers. For a large number of independent identically distributed random variables with finite variance, the average approximately has a normal distribution, no matter what the distribution of the is, with the approximation roughly improving in proportion to .
Example
Suppose is an independent sample from a normally distributed population with unknown parameters mean and variance Let
Where is the sample mean, and is the sample variance. Then
has a Student's t distribution with degrees of freedom. Note that the distribution of does not depend on the values of the unobservable parameters and ; i.e., it is a pivotal quantity. Suppose we wanted to calculate a 95% confidence interval for Then, denoting as the 97.5th percentile of this distribution,
Note that "97.5th" and "0.95" are correct in the preceding expressions. There is a 2.5% chance that will be less than and a 2.5% chance that it will be larger than Thus, the probability that will be between and is 95%. is the probability measure under the student distribution.
Consequently,
and we have a theoretical (stochastic) 95% confidence interval for Here is the probability measure under unknown distribution of .
After observing the sample we find values for and for from which we compute the confidence interval
Interpretation
Various interpretations of a confidence interval can be given (taking the 95% confidence interval as an example in the following).
The confidence interval can be expressed in terms of a long-run frequency in repeated samples (or in resampling): "Were this procedure to be repeated on numerous samples, the proportion of calculated 95% confidence intervals that encompassed the true value of the population parameter would tend toward 95%."
The confidence interval can be expressed in terms of probability with respect to a single theoretical (yet to be realized) sample: "There is a 95% probability that the 95% confidence interval calculated from a given future sample will cover the true value of the population parameter." This essentially reframes the "repeated samples" interpretation as a probability rather than a frequency.
The confidence interval can be expressed in terms of statistical significance, e.g.: "The 95% confidence interval represents values that are not statistically significantly different from the point estimate at the .05 level."
Common misunderstandings
Confidence intervals and levels are frequently misunderstood, and published studies have shown that even professional scientists often misinterpret them.
A 95% confidence level does not mean that for a given realized interval there is a 95% probability that the population parameter lies within the interval (i.e., a 95% probability that the interval covers the population parameter). According to the frequentist interpretation, once an interval is calculated, this interval either covers the parameter value or it does not; it is no longer a matter of probability. The 95% probability relates to the reliability of the estimation procedure, not to a specific calculated interval. Neyman himself (the original proponent of confidence intervals) made this point in his original paper:It will be noticed that in the above description, the probability statements refer to the problems of estimation with which the statistician will be concerned in the future. In fact, I have repeatedly stated that the frequency of correct results will tend to α. Consider now the case when a sample is already drawn, and the calculations have given [particular limits]. Can we say that in this particular case the probability of the true value [falling between these limits] is equal to α? The answer is obviously in the negative. The parameter is an unknown constant, and no probability statement concerning its value may be made...
A 95% confidence level does not mean that 95% of the sample data lie within the confidence interval.
A 95% confidence level does not mean that there is a 95% probability of the parameter estimate from a repeat of the experiment falling within the confidence interval computed from a given experiment.
Examples of how naïve interpretation of confidence intervals can be problematic
Confidence procedure for uniform location
Welch presented an example which clearly shows the difference between the theory of confidence intervals and other theories of interval estimation (including Fisher's fiducial intervals and objective Bayesian intervals). Robinson called this example "[p]ossibly the best known counterexample for Neyman's version of confidence interval theory." To Welch, it showed the superiority of confidence interval theory; to critics of the theory, it shows a deficiency. Here we present a simplified version.
Suppose that are independent observations from a uniform distribution. Then the optimal 50% confidence procedure for is
A fiducial or objective Bayesian argument can be used to derive the interval estimate
which is also a 50% confidence procedure. Welch showed that the first confidence procedure dominates the second, according to desiderata from confidence interval theory; for every , the probability that the first procedure contains is less than or equal to the probability that the second procedure contains . The average width of the intervals from the first procedure is less than that of the second. Hence, the first procedure is preferred under classical confidence interval theory.
However, when , intervals from the first procedure are guaranteed to contain the true value : Therefore, the nominal 50% confidence coefficient is unrelated to the uncertainty we should have that a specific interval contains the true value. The second procedure does not have this property.
Moreover, when the first procedure generates a very short interval, this indicates that are very close together and hence only offer the information in a single data point. Yet the first interval will exclude almost all reasonable values of the parameter due to its short width. The second procedure does not have this property.
The two counter-intuitive properties of the first procedure – 100% coverage when are far apart and almost 0% coverage when are close together – balance out to yield 50% coverage on average. However, despite the first procedure being optimal, its intervals offer neither an assessment of the precision of the estimate nor an assessment of the uncertainty one should have that the interval contains the true value.
This example is used to argue against naïve interpretations of confidence intervals. If a confidence procedure is asserted to have properties beyond that of the nominal coverage (such as relation to precision, or a relationship with Bayesian inference), those properties must be proved; they do not follow from the fact that a procedure is a confidence procedure.
Confidence procedure for ω2
Steiger suggested a number of confidence procedures for common effect size measures in ANOVA. Morey et al. point out that several of these confidence procedures, including the one for ω2, have the property that as the F statistic becomes increasingly small—indicating misfit with all possible values of ω2—the confidence interval shrinks and can even contain only the single value ω2 = 0; that is, the CI is infinitesimally narrow (this occurs when for a CI).
This behavior is consistent with the relationship between the confidence procedure and significance testing: as F becomes so small that the group means are much closer together than we would expect by chance, a significance test might indicate rejection for most or all values of ω2. Hence the interval will be very narrow or even empty (or, by a convention suggested by Steiger, containing only 0). However, this does not indicate that the estimate of ω2 is very precise. In a sense, it indicates the opposite: that the trustworthiness of the results themselves may be in doubt. This is contrary to the common interpretation of confidence intervals that they reveal the precision of the estimate.
Confidence interval for specific distributions
Confidence interval for binomial distribution
Confidence interval for exponent of the power law distribution
Confidence interval for mean of the exponential distribution
Confidence interval for mean of the Poisson distribution
Confidence intervals for mean and variance of the normal distribution (also here)
Confidence interval for the parameters of a simple linear regression
Confidence interval for the difference of means (based on data from a normal distributions, without assuming equal variances)
Comparing the Proportions of Two Binomials using z-test
| Mathematics | Statistics | null |
19826952 | https://en.wikipedia.org/wiki/Protostome | Protostome | Protostomia () is the clade of animals once thought to be characterized by the formation of the organism's mouth before its anus during embryonic development. This nature has since been discovered to be extremely variable among Protostomia's members, although the reverse is typically true of its sister clade, Deuterostomia. Well-known examples of protostomes are arthropods, molluscs, annelids, flatworms and nematodes. They are also called schizocoelomates since schizocoely typically occurs in them.
Together with the Deuterostomia and Xenacoelomorpha, these form the clade Bilateria, animals with bilateral symmetry, anteroposterior axis and three germ layers.
Protostomy
In animals at least as complex as earthworms, the first phase in gut development involves the embryo forming a dent on one side (the blastopore) which deepens to become its digestive tube (the archenteron). In the sister-clade, the deuterostomes (), the original dent becomes the anus while the gut eventually tunnels through to make another opening, which forms the mouth. The protostomes (from Greek 'first' + 'mouth') were so named because it was once believed that in all cases the embryological dent formed the mouth while the anus was formed later, at the opening made by the other end of the gut.
It is now known that the fate of the blastopore among protostomes is extremely variable; while the evolutionary distinction between deuterostomes and protostomes remains valid, the descriptive accuracy of the name protostome is disputable.
Protostome and deuterostome embryos differ in several other ways. Secondary body cavities (coeloms) generally form by schizocoely, where the coelom forms out of a solid mass of embryonic tissue splitting away from the rest, instead of by enterocoelic pouching, where the coelom would otherwise form out of in-folded gut walls.
Evolution
The common ancestor of protostomes and deuterostomes was evidently a worm-like aquatic animal of the Ediacaran. The two clades diverged about 600 million years ago. Protostomes evolved into over a million species alive today, compared to ca. 73,000 deuterostome species.
Protostomes are divided into the Ecdysozoa (e.g. arthropods, nematodes) and the Spiralia (e.g. molluscs, annelids, platyhelminths, and rotifers). A modern consensus phylogenetic tree for the protostomes is shown below.
The timing of clades radiating into newer clades is given in mya (millions of years ago); less certain placements are indicated with dashed lines.
| Biology and health sciences | General classifications | Animals |
19827221 | https://en.wikipedia.org/wiki/Arthropod | Arthropod | Arthropods ( ) are invertebrates in the phylum Arthropoda. They possess an exoskeleton with a cuticle made of chitin, often mineralised with calcium carbonate, a body with differentiated (metameric) segments, and paired jointed appendages. In order to keep growing, they must go through stages of moulting, a process by which they shed their exoskeleton to reveal a new one. They form an extremely diverse group of up to ten million species.
Haemolymph is the analogue of blood for most arthropods. An arthropod has an open circulatory system, with a body cavity called a haemocoel through which haemolymph circulates to the interior organs. Like their exteriors, the internal organs of arthropods are generally built of repeated segments. They have ladder-like nervous systems, with paired ventral nerve cords running through all segments and forming paired ganglia in each segment. Their heads are formed by fusion of varying numbers of segments, and their brains are formed by fusion of the ganglia of these segments and encircle the esophagus. The respiratory and excretory systems of arthropods vary, depending as much on their environment as on the subphylum to which they belong.
Arthropods use combinations of compound eyes and pigment-pit ocelli for vision. In most species, the ocelli can only detect the direction from which light is coming, and the compound eyes are the main source of information, but the main eyes of spiders are ocelli that can form images and, in a few cases, can swivel to track prey. Arthropods also have a wide range of chemical and mechanical sensors, mostly based on modifications of the many bristles known as setae that project through their cuticles. Similarly, their reproduction and development are varied; all terrestrial species use internal fertilization, but this is sometimes by indirect transfer of the sperm via an appendage or the ground, rather than by direct injection. Aquatic species use either internal or external fertilization. Almost all arthropods lay eggs, with many species giving birth to live young after the eggs have hatched inside the mother; but a few are genuinely viviparous, such as aphids. Arthropod hatchlings vary from miniature adults to grubs and caterpillars that lack jointed limbs and eventually undergo a total metamorphosis to produce the adult form. The level of maternal care for hatchlings varies from nonexistent to the prolonged care provided by social insects.
The evolutionary ancestry of arthropods dates back to the Cambrian period. The group is generally regarded as monophyletic, and many analyses support the placement of arthropods with cycloneuralians (or their constituent clades) in a superphylum Ecdysozoa. Overall, however, the basal relationships of animals are not yet well resolved. Likewise, the relationships between various arthropod groups are still actively debated. Today, arthropods contribute to the human food supply both directly as food, and more importantly, indirectly as pollinators of crops. Some species are known to spread severe disease to humans, livestock, and crops.
Etymology
The word arthropod comes from the Greek , and (gen. ) or , which together mean "jointed leg", with the word "arthropodes" initially used in anatomical descriptions by Barthélemy Charles Joseph Dumortier published in 1832. The designation "Arthropoda" appears to have been first used in 1843 by the German zoologist Johann Ludwig Christian Gravenhorst (1777–1857). The origin of the name has been the subject of considerable confusion, with credit often given erroneously to Pierre André Latreille or Karl Theodor Ernst von Siebold instead, among various others.
Terrestrial arthropods are often called bugs. The term is also occasionally extended to colloquial names for freshwater or marine crustaceans (e.g., Balmain bug, Moreton Bay bug, mudbug) and used by physicians and bacteriologists for disease-causing germs (e.g., superbugs), but entomologists reserve this term for a narrow category of "true bugs", insects of the order Hemiptera.
Description
Arthropods are invertebrates with segmented bodies and jointed limbs. The exoskeleton or cuticles consists of chitin, a polymer of N-Acetylglucosamine. The cuticle of many crustaceans, beetle mites, the clades Penetini and Archaeoglenini inside the beetle subfamily Phrenapatinae, and millipedes (except for bristly millipedes) is also biomineralized with calcium carbonate. Calcification of the endosternite, an internal structure used for muscle attachments, also occur in some opiliones, and the pupal cuticle of the fly Bactrocera dorsalis contains calcium phosphate.
Diversity
Arthropoda is the largest animal phylum with the estimates of the number of arthropod species varying from 1,170,000 to 5~10 million and accounting for over 80 percent of all known living animal species. One arthropod sub-group, the insects, includes more described species than any other taxonomic class. The total number of species remains difficult to determine. This is due to the census modeling assumptions projected onto other regions in order to scale up from counts at specific locations applied to the whole world. A study in 1992 estimated that there were 500,000 species of animals and plants in Costa Rica alone, of which 365,000 were arthropods.
They are important members of marine, freshwater, land and air ecosystems and one of only two major animal groups that have adapted to life in dry environments; the other is amniotes, whose living members are reptiles, birds and mammals. Both the smallest and largest arthropods are crustaceans. The smallest belong to the class Tantulocarida, some of which are less than long. The largest are species in the class Malacostraca, with the legs of the Japanese spider crab potentially spanning up to and the American lobster reaching weights over 20 kg (44 lbs).
Segmentation
The embryos of all arthropods are segmented, built from a series of repeated modules. The last common ancestor of living arthropods probably consisted of a series of undifferentiated segments, each with a pair of appendages that functioned as limbs. However, all known living and fossil arthropods have grouped segments into tagmata in which segments and their limbs are specialized in various ways.
The three-part appearance of many insect bodies and the two-part appearance of spiders is a result of this grouping. There are no external signs of segmentation in mites. Arthropods also have two body elements that are not part of this serially repeated pattern of segments, an ocular somite at the front, where the mouth and eyes originated, and a telson at the rear, behind the anus.
Originally, it seems that each appendage-bearing segment had two separate pairs of appendages: an upper, unsegmented exite and a lower, segmented endopod. These would later fuse into a single pair of biramous appendages united by a basal segment (protopod or basipod), with the upper branch acting as a gill while the lower branch was used for locomotion. The appendages of most crustaceans and some extinct taxa such as trilobites have another segmented branch known as exopods, but whether these structures have a single origin remain controversial. In some segments of all known arthropods, the appendages have been modified, for example to form gills, mouth-parts, antennae for collecting information, or claws for grasping; arthropods are "like Swiss Army knives, each equipped with a unique set of specialized tools." In many arthropods, appendages have vanished from some regions of the body; it is particularly common for abdominal appendages to have disappeared or be highly modified.
The most conspicuous specialization of segments is in the head. The four major groups of arthropods – Chelicerata (sea spiders, horseshoe crabs and arachnids), Myriapoda (symphylans, pauropods, millipedes and centipedes), Pancrustacea (oligostracans, copepods, malacostracans, branchiopods, hexapods, etc.), and the extinct Trilobita – have heads formed of various combinations of segments, with appendages that are missing or specialized in different ways. Despite myriapods and hexapods both having similar head combinations, hexapods are deeply nested within crustacea while myriapods are not, so these traits are believed to have evolved separately. In addition, some extinct arthropods, such as Marrella, belong to none of these groups, as their heads are formed by their own particular combinations of segments and specialized appendages.
Working out the evolutionary stages by which all these different combinations could have appeared is so difficult that it has long been known as "The arthropod head problem". In 1960, R. E. Snodgrass even hoped it would not be solved, as he found trying to work out solutions to be fun.
Exoskeleton
Arthropod exoskeletons are made of cuticle, a non-cellular material secreted by the epidermis. Their cuticles vary in the details of their structure, but generally consist of three main layers: the epicuticle, a thin outer waxy coat that moisture-proofs the other layers and gives them some protection; the exocuticle, which consists of chitin and chemically hardened proteins; and the endocuticle, which consists of chitin and unhardened proteins. The exocuticle and endocuticle together are known as the procuticle. Each body segment and limb section is encased in hardened cuticle. The joints between body segments and between limb sections are covered by flexible cuticle.
The exoskeletons of most aquatic crustaceans are biomineralized with calcium carbonate extracted from the water. Some terrestrial crustaceans have developed means of storing the mineral, since on land they cannot rely on a steady supply of dissolved calcium carbonate. Biomineralization generally affects the exocuticle and the outer part of the endocuticle. Two recent hypotheses about the evolution of biomineralization in arthropods and other groups of animals propose that it provides tougher defensive armor, and that it allows animals to grow larger and stronger by providing more rigid skeletons; and in either case a mineral-organic composite exoskeleton is cheaper to build than an all-organic one of comparable strength.
The cuticle may have setae (bristles) growing from special cells in the epidermis. Setae are as varied in form and function as appendages. For example, they are often used as sensors to detect air or water currents, or contact with objects; aquatic arthropods use feather-like setae to increase the surface area of swimming appendages and to filter food particles out of water; aquatic insects, which are air-breathers, use thick felt-like coats of setae to trap air, extending the time they can spend under water; heavy, rigid setae serve as defensive spines.
Although all arthropods use muscles attached to the inside of the exoskeleton to flex their limbs, some still use hydraulic pressure to extend them, a system inherited from their pre-arthropod ancestors; for example, all spiders extend their legs hydraulically and can generate pressures up to eight times their resting level.
Moulting
The exoskeleton cannot stretch and thus restricts growth. Arthropods, therefore, replace their exoskeletons by undergoing ecdysis (moulting), or shedding the old exoskeleton, the exuviae, after growing a new one that is not yet hardened. Moulting cycles run nearly continuously until an arthropod reaches full size. The developmental stages between each moult (ecdysis) until sexual maturity is reached is called an instar. Differences between instars can often be seen in altered body proportions, colors, patterns, changes in the number of body segments or head width. After moulting, i.e. shedding their exoskeleton, the juvenile arthropods continue in their life cycle until they either pupate or moult again.
In the initial phase of moulting, the animal stops feeding and its epidermis releases moulting fluid, a mixture of enzymes that digests the endocuticle and thus detaches the old cuticle. This phase begins when the epidermis has secreted a new epicuticle to protect it from the enzymes, and the epidermis secretes the new exocuticle while the old cuticle is detaching. When this stage is complete, the animal makes its body swell by taking in a large quantity of water or air, and this makes the old cuticle split along predefined weaknesses where the old exocuticle was thinnest. It commonly takes several minutes for the animal to struggle out of the old cuticle. At this point, the new one is wrinkled and so soft that the animal cannot support itself and finds it very difficult to move, and the new endocuticle has not yet formed. The animal continues to pump itself up to stretch the new cuticle as much as possible, then hardens the new exocuticle and eliminates the excess air or water. By the end of this phase, the new endocuticle has formed. Many arthropods then eat the discarded cuticle to reclaim its materials.
Because arthropods are unprotected and nearly immobilized until the new cuticle has hardened, they are in danger both of being trapped in the old cuticle and of being attacked by predators. Moulting may be responsible for 80 to 90% of all arthropod deaths.
Internal organs
Arthropod bodies are also segmented internally, and the nervous, muscular, circulatory, and excretory systems have repeated components. Arthropods come from a lineage of animals that have a coelom, a membrane-lined cavity between the gut and the body wall that accommodates the internal organs. The strong, segmented limbs of arthropods eliminate the need for one of the coelom's main ancestral functions, as a hydrostatic skeleton, which muscles compress in order to change the animal's shape and thus enable it to move. Hence the coelom of the arthropod is reduced to small areas around the reproductive and excretory systems. Its place is largely taken by a hemocoel, a cavity that runs most of the length of the body and through which blood flows.
Respiration and circulation
Arthropods have open circulatory systems. Most have a few short, open-ended arteries. In chelicerates and crustaceans, the blood carries oxygen to the tissues, while hexapods use a separate system of tracheae. Many crustaceans and a few chelicerates and tracheates use respiratory pigments to assist oxygen transport. The most common respiratory pigment in arthropods is copper-based hemocyanin; this is used by many crustaceans and a few centipedes. A few crustaceans and insects use iron-based hemoglobin, the respiratory pigment used by vertebrates. As with other invertebrates, the respiratory pigments of those arthropods that have them are generally dissolved in the blood and rarely enclosed in corpuscles as they are in vertebrates.
The heart is a muscular tube that runs just under the back and for most of the length of the hemocoel. It contracts in ripples that run from rear to front, pushing blood forwards. Sections not being squeezed by the heart muscle are expanded either by elastic ligaments or by small muscles, in either case connecting the heart to the body wall. Along the heart run a series of paired ostia, non-return valves that allow blood to enter the heart but prevent it from leaving before it reaches the front.
Arthropods have a wide variety of respiratory systems. Small species often do not have any, since their high ratio of surface area to volume enables simple diffusion through the body surface to supply enough oxygen. Crustacea usually have gills that are modified appendages. Many arachnids have book lungs. Tracheae, systems of branching tunnels that run from the openings in the body walls, deliver oxygen directly to individual cells in many insects, myriapods and arachnids.
Nervous system
Living arthropods have paired main nerve cords running along their bodies below the gut, and in each segment the cords form a pair of ganglia from which sensory and motor nerves run to other parts of the segment. Although the pairs of ganglia in each segment often appear physically fused, they are connected by commissures (relatively large bundles of nerves), which give arthropod nervous systems a characteristic ladder-like appearance. The brain is in the head, encircling and mainly above the esophagus. It consists of the fused ganglia of the acron and one or two of the foremost segments that form the head – a total of three pairs of ganglia in most arthropods, but only two in chelicerates, which do not have antennae or the ganglion connected to them. The ganglia of other head segments are often close to the brain and function as part of it. In insects, these other head ganglia combine into a pair of subesophageal ganglia, under and behind the esophagus. Spiders take this process a step further, as all the segmental ganglia are incorporated into the subesophageal ganglia, which occupy most of the space in the cephalothorax (front "super-segment").
Excretory system
There are two different types of arthropod excretory systems. In aquatic arthropods, the end-product of biochemical reactions that metabolise nitrogen is ammonia, which is so toxic that it needs to be diluted as much as possible with water. The ammonia is then eliminated via any permeable membrane, mainly through the gills. All crustaceans use this system, and its high consumption of water may be responsible for the relative lack of success of crustaceans as land animals. Various groups of terrestrial arthropods have independently developed a different system: the end-product of nitrogen metabolism is uric acid, which can be excreted as dry material; the Malpighian tubule system filters the uric acid and other nitrogenous waste out of the blood in the hemocoel, and dumps these materials into the hindgut, from which they are expelled as feces. Most aquatic arthropods and some terrestrial ones also have organs called nephridia ("little kidneys"), which extract other wastes for excretion as urine.
Senses
The stiff cuticles of arthropods would block out information about the outside world, except that they are penetrated by many sensors or connections from sensors to the nervous system. In fact, arthropods have modified their cuticles into elaborate arrays of sensors. Various touch sensors, mostly setae, respond to different levels of force, from strong contact to very weak air currents. Chemical sensors provide equivalents of taste and smell, often by means of setae. Pressure sensors often take the form of membranes that function as eardrums, but are connected directly to nerves rather than to auditory ossicles. The antennae of most hexapods include sensor packages that monitor humidity, moisture and temperature.
Most arthropods lack balance and acceleration sensors, and rely on their eyes to tell them which way is up. The self-righting behavior of cockroaches is triggered when pressure sensors on the underside of the feet report no pressure. However, many malacostracan crustaceans have statocysts, which provide the same sort of information as the balance and motion sensors of the vertebrate inner ear.
The proprioceptors of arthropods, sensors that report the force exerted by muscles and the degree of bending in the body and joints, are well understood. However, little is known about what other internal sensors arthropods may have.
Optical
Most arthropods have sophisticated visual systems that include one or more usually both of compound eyes and pigment-cup ocelli ("little eyes"). In most cases, ocelli are only capable of detecting the direction from which light is coming, using the shadow cast by the walls of the cup. However, the main eyes of spiders are pigment-cup ocelli that are capable of forming images, and those of jumping spiders can rotate to track prey.
Compound eyes consist of fifteen to several thousand independent ommatidia, columns that are usually hexagonal in cross section. Each ommatidium is an independent sensor, with its own light-sensitive cells and often with its own lens and cornea. Compound eyes have a wide field of view, and can detect fast movement and, in some cases, the polarization of light. On the other hand, the relatively large size of ommatidia makes the images rather coarse, and compound eyes are shorter-sighted than those of birds and mammals – although this is not a severe disadvantage, as objects and events within are most important to most arthropods. Several arthropods have color vision, and that of some insects has been studied in detail; for example, the ommatidia of bees contain receptors for both green and ultra-violet.
Olfaction
Reproduction and development
A few arthropods, such as barnacles, are hermaphroditic, that is, each can have the organs of both sexes. However, individuals of most species remain of one sex their entire lives. A few species of insects and crustaceans can reproduce by parthenogenesis, especially if conditions favor a "population explosion". However, most arthropods rely on sexual reproduction, and parthenogenetic species often revert to sexual reproduction when conditions become less favorable. The ability to undergo meiosis is widespread among arthropods including both those that reproduce sexually and those that reproduce parthenogenetically. Although meiosis is a major characteristic of arthropods, understanding of its fundamental adaptive benefit has long been regarded as an unresolved problem, that appears to have remained unsettled.
Aquatic arthropods may breed by external fertilization, as for example horseshoe crabs do, or by internal fertilization, where the ova remain in the female's body and the sperm must somehow be inserted. All known terrestrial arthropods use internal fertilization. Opiliones (harvestmen), millipedes, and some crustaceans use modified appendages such as gonopods or penises to transfer the sperm directly to the female. However, most male terrestrial arthropods produce spermatophores, waterproof packets of sperm, which the females take into their bodies. A few such species rely on females to find spermatophores that have already been deposited on the ground, but in most cases males only deposit spermatophores when complex courtship rituals look likely to be successful.
Most arthropods lay eggs, but scorpions are ovoviviparous: they produce live young after the eggs have hatched inside the mother, and are noted for prolonged maternal care. Newly born arthropods have diverse forms, and insects alone cover the range of extremes. Some hatch as apparently miniature adults (direct development), and in some cases, such as silverfish, the hatchlings do not feed and may be helpless until after their first moult. Many insects hatch as grubs or caterpillars, which do not have segmented limbs or hardened cuticles, and metamorphose into adult forms by entering an inactive phase in which the larval tissues are broken down and re-used to build the adult body. Dragonfly larvae have the typical cuticles and jointed limbs of arthropods but are flightless water-breathers with extendable jaws. Crustaceans commonly hatch as tiny nauplius larvae that have only three segments and pairs of appendages.
Evolutionary history
Last common ancestor
Based on the distribution of shared plesiomorphic features in extant and fossil taxa, the last common ancestor of all arthropods is inferred to have been as a modular organism with each module covered by its own sclerite (armor plate) and bearing a pair of biramous limbs. However, whether the ancestral limb was uniramous or biramous is far from a settled debate.
This Ur-arthropod had a ventral mouth, pre-oral antennae and dorsal eyes at the front of the body. It was assumed to have been a non-discriminatory sediment feeder, processing whatever sediment came its way for food, but fossil findings hint that the last common ancestor of both arthropods and Priapulida shared the same specialized mouth apparatus: a circular mouth with rings of teeth used for capturing animal prey.
Fossil record
It has been proposed that the Ediacaran animals Parvancorina and Spriggina, from around , were arthropods, but later study shows that their affinities of being origin of arthropods are not reliable. Small arthropods with bivalve-like shells have been found in Early Cambrian fossil beds dating in China and Australia. The earliest Cambrian trilobite fossils are about 520 million years old, but the class was already quite diverse and worldwide, suggesting that they had been around for quite some time. In the Maotianshan shales, which date back to 518 million years ago, arthropods such as Kylinxia and Erratus have been found that seem to represent transitional fossils between stem (e.g. Radiodonta such as Anomalocaris) and true arthropods. Re-examination in the 1970s of the Burgess Shale fossils from about identified many arthropods, some of which could not be assigned to any of the well-known groups, and thus intensified the debate about the Cambrian explosion. A fossil of Marrella from the Burgess Shale has provided the earliest clear evidence of moulting.
The earliest fossil of likely pancrustacean larvae date from about in the Cambrian, followed by unique taxa like Yicaris and Wujicaris. The purported pancrustacean/crustacean affinity of some cambrian arthropods (e.g. Phosphatocopina, Bradoriida and Hymenocarine taxa like waptiids) were disputed by subsequent studies, as they might branch before the mandibulate crown-group. Within the pancrustacean crown-group, only Malacostraca, Branchiopoda and Pentastomida have Cambrian fossil records. Crustacean fossils are common from the Ordovician period onwards. They have remained almost entirely aquatic, possibly because they never developed excretory systems that conserve water.
Arthropods provide the earliest identifiable fossils of land animals, from about in the Late Silurian, and terrestrial tracks from about appear to have been made by arthropods. Arthropods possessed attributes that were easy coopted for life on land; their existing jointed exoskeletons provided protection against desiccation, support against gravity and a means of locomotion that was not dependent on water. Around the same time the aquatic, scorpion-like eurypterids became the largest ever arthropods, some as long as .
The oldest known arachnid is the trigonotarbid Palaeotarbus jerami, from about in the Silurian period. Attercopus fimbriunguis, from in the Devonian period, bears the earliest known silk-producing spigots, but its lack of spinnerets means it was not one of the true spiders, which first appear in the Late Carboniferous over . The Jurassic and Cretaceous periods provide a large number of fossil spiders, including representatives of many modern families. The oldest known scorpion is Dolichophonus, dated back to . Lots of Silurian and Devonian scorpions were previously thought to be gill-breathing, hence the idea that scorpions were primitively aquatic and evolved air-breathing book lungs later on. However subsequent studies reveal most of them lacking reliable evidence for an aquatic lifestyle, while exceptional aquatic taxa (e.g. Waeringoscorpio) most likely derived from terrestrial scorpion ancestors.
The oldest fossil record of hexapod is obscure, as most of the candidates are poorly preserved and their hexapod affinities had been disputed. An iconic example is the Devonian Rhyniognatha hirsti, dated at , its mandibles are thought to be a type found only in winged insects, which suggests that the earliest insects appeared in the Silurian period. However later study shows that Rhyniognatha most likely represent a myriapod, not even a hexapod. The unequivocal oldest known hexapod is the springtail Rhyniella, from about in the Devonian period, and the palaeodictyopteran Delitzschala bitterfeldensis, from about in the Carboniferous period, respectively. The Mazon Creek lagerstätten from the Late Carboniferous, about , include about 200 species, some gigantic by modern standards, and indicate that insects had occupied their main modern ecological niches as herbivores, detritivores and insectivores. Social termites and ants first appear in the Early Cretaceous, and advanced social bees have been found in Late Cretaceous rocks but did not become abundant until the Middle Cenozoic.
Evolutionary relationships to other animal phyla
From 1952 to 1977, zoologist Sidnie Manton and others argued that arthropods are polyphyletic, in other words, that they do not share a common ancestor that was itself an arthropod. Instead, they proposed that three separate groups of "arthropods" evolved separately from common worm-like ancestors: the chelicerates, including spiders and scorpions; the crustaceans; and the uniramia, consisting of onychophorans, myriapods and hexapods. These arguments usually bypassed trilobites, as the evolutionary relationships of this class were unclear. Proponents of polyphyly argued the following: that the similarities between these groups are the results of convergent evolution, as natural consequences of having rigid, segmented exoskeletons; that the three groups use different chemical means of hardening the cuticle; that there were significant differences in the construction of their compound eyes; that it is hard to see how such different configurations of segments and appendages in the head could have evolved from the same ancestor; and that crustaceans have biramous limbs with separate gill and leg branches, while the other two groups have uniramous limbs in which the single branch serves as a leg.
Further analysis and discoveries in the 1990s reversed this view, and led to acceptance that arthropods are monophyletic, in other words they are inferred to share a common ancestor that was itself an arthropod. For example, Graham Budd's analyses of Kerygmachela in 1993 and of Opabinia in 1996 convinced him that these animals were similar to onychophorans and to various Early Cambrian "lobopods", and he presented an "evolutionary family tree" that showed these as "aunts" and "cousins" of all arthropods. These changes made the scope of the term "arthropod" unclear, and Claus Nielsen proposed that the wider group should be labelled "Panarthropoda" ("all the arthropods") while the animals with jointed limbs and hardened cuticles should be called "Euarthropoda" ("true arthropods").
A contrary view was presented in 2003, when Jan Bergström and Hou Xian-guang argued that, if arthropods were a "sister-group" to any of the anomalocarids, they must have lost and then re-evolved features that were well-developed in the anomalocarids. The earliest known arthropods ate mud in order to extract food particles from it, and possessed variable numbers of segments with unspecialized appendages that functioned as both gills and legs. Anomalocarids were, by the standards of the time, huge and sophisticated predators with specialized mouths and grasping appendages, fixed numbers of segments some of which were specialized, tail fins, and gills that were very different from those of arthropods. In 2006, they suggested that arthropods were more closely related to lobopods and tardigrades than to anomalocarids. In 2014, it was found that tardigrades were more closely related to arthropods than velvet worms.
Higher up the "family tree", the Annelida have traditionally been considered the closest relatives of the Panarthropoda, since both groups have segmented bodies, and the combination of these groups was labelled Articulata. There had been competing proposals that arthropods were closely related to other groups such as nematodes, priapulids and tardigrades, but these remained minority views because it was difficult to specify in detail the relationships between these groups.
In the 1990s, molecular phylogenetic analyses of DNA sequences produced a coherent scheme showing arthropods as members of a superphylum labelled Ecdysozoa ("animals that moult"), which contained nematodes, priapulids and tardigrades but excluded annelids. This was backed up by studies of the anatomy and development of these animals, which showed that many of the features that supported the Articulata hypothesis showed significant differences between annelids and the earliest Panarthropods in their details, and some were hardly present at all in arthropods. This hypothesis groups annelids with molluscs and brachiopods in another superphylum, Lophotrochozoa.
If the Ecdysozoa hypothesis is correct, then segmentation of arthropods and annelids either has evolved convergently or has been inherited from a much older ancestor and subsequently lost in several other lineages, such as the non-arthropod members of the Ecdysozoa.
Evolution of fossil arthropods
Aside from the four major living groups (crustaceans, chelicerates, myriapods and hexapods), a number of fossil forms, mostly from the early Cambrian period, are difficult to place taxonomically, either from lack of obvious affinity to any of the main groups or from clear affinity to several of them. Marrella was the first one to be recognized as significantly different from the well-known groups.
Modern interpretations of the basal, extinct stem-group of Arthropoda recognised the following groups, from most basal to most crownward:
The "Giant" or "Siberiid Lobopodians", such as Jianshanopodia, Siberion and Megadictyon, are the most basal grade in the total-group Arthropoda.
The "Gilled Lobopodians", such as Kerygmachela, Pambdelurion and Opabinia, are the second most basal grade.
The Radiodonta, which traditionally known as anomalocaridids come in third position, and are thought to be monophyletic.
A possible "upper stem-group" assemblage of more uncertain position but contained within Deuteropoda: the Fuxianhuiida, Megacheira, and multiple "bivalved forms" including Isoxyida and Hymenocarina.
The Deuteropoda is a recently established clade uniting the crown-group (living) arthropods with these possible "upper stem-group" fossils taxa. The clade is defined by important changes to the structure of the head region such as the appearance of a differentiated deutocerebral appendage pair, which excludes more basal taxa like radiodonts and "gilled lobopodians".
Controversies remain about the positions of various extinct arthropod groups. Some studies recover Megacheira as closely related to chelicerates, while others recover them as outside the group containing Chelicerate and Mandibulata as stem-group euarthropods. The placement of the Artiopoda (which contains the extinct trilobites and similar forms) is also a frequent subject of dispute. The main hypotheses position them in the clade Arachnomorpha with the Chelicerates. However, one of the newer hypotheses is that the chelicerae have originated from the same pair of appendages that evolved into antennae in the ancestors of Mandibulata, which would place trilobites, which had antennae, closer to Mandibulata than Chelicerata, in the clade Antennulata. The fuxianhuiids, usually suggested to be stem-group arthropods, have been suggested to be Mandibulates in some recent studies. The Hymenocarina, a group of bivalved arthropods, previously thought to have been stem-group members of the group, have been demonstrated to be mandibulates based on the presence of mandibles.
Evolution and classification of living arthropods
The phylum Arthropoda is typically subdivided into four subphyla, of which one is extinct:
Artiopods are an extinct group of formerly numerous marine arthropods that disappeared in the Permian–Triassic extinction event, though they were in decline prior to this killing blow, having been reduced to a handful of orders in the Late Devonian extinction. They contain groups such as the trilobites, nektaspids, aglaspidids, and the cheloniellids among others.
Chelicerates comprise the marine sea spiders and horseshoe crabs, along with the terrestrial arachnids such as mites, harvestmen, spiders, scorpions and related organisms characterized by the presence of chelicerae, appendages just above/in front of the mouthparts. Chelicerae appear in scorpions and horseshoe crabs as tiny claws that they use in feeding, but those of spiders have developed as fangs that inject venom.
Myriapods comprise millipedes, centipedes, pauropods and symphylans, characterized by having numerous body segments each of which bearing one or two pairs of legs (or in a few cases being legless). All members are exclusively terrestrial.
Pancrustaceans comprise ostracods, barnacles, copepods, malacostracans, cephalocaridans, branchiopods, remipedes and hexapods. Most groups are primarily aquatic (two notable exceptions being woodlice and hexapods, which are both purely terrestrial) and are characterized by having biramous appendages. The most abundant group of pancrustaceans are the terrestrial hexapods, which comprise insects, diplurans, springtails, and proturans, with six thoracic legs.
The phylogeny of the major extant arthropod groups has been an area of considerable interest and dispute. Recent studies strongly suggest that Crustacea, as traditionally defined, is paraphyletic, with Hexapoda having evolved from within it, so that Crustacea and Hexapoda form a clade, Pancrustacea. The position of Myriapoda, Chelicerata and Pancrustacea remains unclear . In some studies, Myriapoda is grouped with Chelicerata (forming Myriochelata); in other studies, Myriapoda is grouped with Pancrustacea (forming Mandibulata), or Myriapoda may be sister to Chelicerata plus Pancrustacea.
The following cladogram shows the internal relationships between all the living classes of arthropods as of the late 2010s, as well as the estimated timing for some of the clades:
Interaction with humans
Crustaceans such as crabs, lobsters, crayfish, shrimp, and prawns have long been part of human cuisine, and are now raised commercially. Insects and their grubs are at least as nutritious as meat, and are eaten both raw and cooked in many cultures, though not most European, Hindu, and Islamic cultures. Cooked tarantulas are considered a delicacy in Cambodia, and by the Piaroa Indians of southern Venezuela, after the highly irritant hairs – the spider's main defense system – are removed. Humans also unintentionally eat arthropods in other foods, and food safety regulations lay down acceptable contamination levels for different kinds of food material. The intentional cultivation of arthropods and other small animals for human food, referred to as minilivestock, is now emerging in animal husbandry as an ecologically sound concept. Commercial butterfly breeding provides Lepidoptera stock to butterfly conservatories, educational exhibits, schools, research facilities, and cultural events.
However, the greatest contribution of arthropods to human food supply is by pollination: a 2008 study examined the 100 crops that FAO lists as grown for food, and estimated pollination's economic value as €153 billion, or 9.5 per cent of the value of world agricultural production used for human food in 2005. Besides pollinating, bees produce honey, which is the basis of a rapidly growing industry and international trade.
The red dye cochineal, produced from a Central American species of insect, was economically important to the Aztecs and Mayans. While the region was under Spanish control, it became Mexico's second most-lucrative export, and is now regaining some of the ground it lost to synthetic competitors. Shellac, a resin secreted by a species of insect native to southern Asia, was historically used in great quantities for many applications in which it has mostly been replaced by synthetic resins, but it is still used in woodworking and as a food additive. The blood of horseshoe crabs contains a clotting agent, Limulus Amebocyte Lysate, which is now used to test that antibiotics and kidney machines are free of dangerous bacteria, and to detect spinal meningitis. Forensic entomology uses evidence provided by arthropods to establish the time and sometimes the place of death of a human, and in some cases the cause. Recently insects have also gained attention as potential sources of drugs and other medicinal substances.
The relative simplicity of the arthropods' body plan, allowing them to move on a variety of surfaces both on land and in water, have made them useful as models for robotics. The redundancy provided by segments allows arthropods and biomimetic robots to move normally even with damaged or lost appendages.
Although arthropods are the most numerous phylum on Earth, and thousands of arthropod species are venomous, they inflict relatively few serious bites and stings on humans. Far more serious are the effects on humans of diseases like malaria carried by blood-sucking insects. Other blood-sucking insects infect livestock with diseases that kill many animals and greatly reduce the usefulness of others. Ticks can cause tick paralysis and several parasite-borne diseases in humans. A few of the closely related mites also infest humans, causing intense itching, and others cause allergic diseases, including hay fever, asthma, and eczema.
Many species of arthropods, principally insects but also mites, are agricultural and forest pests. The mite Varroa destructor has become the largest single problem faced by beekeepers worldwide. Efforts to control arthropod pests by large-scale use of pesticides have caused long-term effects on human health and on biodiversity. Increasing arthropod resistance to pesticides has led to the development of integrated pest management using a wide range of measures including biological control. Predatory mites may be useful in controlling some mite pests.
| Biology and health sciences | Biology | null |
19827372 | https://en.wikipedia.org/wiki/Brachiopod | Brachiopod | Brachiopods (), phylum Brachiopoda, are a phylum of animals that have hard "valves" (shells) on the upper and lower surfaces, unlike the left and right arrangement in bivalve molluscs. Brachiopod valves are hinged at the rear end, while the front can be opened for feeding or closed for protection. Two major categories are traditionally recognized, articulate and inarticulate brachiopods. The word "articulate" is used to describe the tooth-and-groove structures of the valve-hinge which is present in the articulate group, and absent from the inarticulate group. This is the leading diagnostic skeletal feature, by which the two main groups can be readily distinguished as fossils. Articulate brachiopods have toothed hinges and simple, vertically oriented opening and closing muscles. Conversely, inarticulate brachiopods have weak, untoothed hinges and a more complex system of vertical and oblique (diagonal) muscles used to keep the two valves aligned. In many brachiopods, a stalk-like pedicle projects from an opening near the hinge of one of the valves, known as the pedicle or ventral valve. The pedicle, when present, keeps the animal anchored to the seabed but clear of sediment which would obstruct the opening.
Brachiopod lifespans range from three to over thirty years. Ripe gametes (ova or sperm) float from the gonads into the main coelom and then exit into the mantle cavity. The larvae of inarticulate brachiopods are miniature adults, with lophophores that enable the larvae to feed and swim for months until the animals become heavy enough to settle to the seabed. The planktonic larvae of articulate species do not resemble the adults, but rather look like blobs with yolk sacs, and remain among the plankton for only a few days before leaving the water column upon metamorphosing.
While traditional classification of brachiopods separate them into distinct inarticulate and articulate groups, two approaches appeared in the 1990s. One approach groups the inarticulate Craniida with articulate brachiopods, since both use layers of calcareous minerals their shell; the other approach considers the Craniida to be a separate third group, as their outer organic layer is distinct from that of both the linguliforms ("typical" inarticulates) and rhynchonelliforms (articulates). However, some taxonomists believe it is premature to suggest higher levels of classification such as order and recommend a bottom-up approach that identifies genera and then groups these into intermediate groups. Traditionally, brachiopods have been regarded as members of, or as a sister group to, the deuterostomes, a superphylum that includes chordates and echinoderms. One type of analysis of the evolutionary relationships of brachiopods has always placed brachiopods as protostomes while another type has split between placing brachiopods among the protostomes or the deuterostomes.
It was suggested in 2003 that brachiopods had evolved from an ancestor similar to Halkieria, a slug-like Cambrian animal with "chain mail" on its back and a shell at the front and rear end; it was thought that the ancestral brachiopod converted its shells into a pair of valves by folding the rear part of its body under its front. However, new fossils found in 2007 and 2008 showed that the "chain mail" of tommotiids formed the tube of a sessile animal; one tommotiid resembled phoronids, which are close relatives or a subgroup of brachiopods, while the other tommotiid bore two symmetrical plates that might be an early form of brachiopod valves. Lineages of brachiopods that have both fossil and extant taxa appeared in the early Cambrian, Ordovician, and Carboniferous periods, respectively. Other lineages have arisen and then become extinct, sometimes during severe mass extinctions. At their peak in the Paleozoic era, the brachiopods were among the most abundant filter-feeders and reef-builders, and occupied other ecological niches, including swimming in the jet-propulsion style of scallops. Brachiopod fossils have been useful indicators of climate changes during the Paleozoic. However, after the Permian–Triassic extinction event, brachiopods recovered only a third of their former diversity. A study in 2007 concluded the brachiopods were especially vulnerable to the Permian–Triassic extinction, as they built calcareous hard parts (made of calcium carbonate) and had low metabolic rates and weak respiratory systems. It was often thought that brachiopods went into decline after the Permian–Triassic extinction, and were out-competed by bivalves, but a study in 1980 found both brachiopod and bivalve species increased from the Paleozoic to modern times, with bivalves increasing faster; after the Permian–Triassic extinction, brachiopods became for the first time less diverse than bivalves.
Brachiopods live only in the sea, and most species avoid locations with strong currents or waves. The larvae of articulate species settle in quickly and form dense populations in well-defined areas while the larvae of inarticulate species swim for up to a month and have wide ranges. Brachiopods now live mainly in cold water and low light. Fish and crustaceans seem to find brachiopod flesh distasteful and seldom attack them. Among brachiopods, only the lingulids (Lingula sp.) have been fished commercially, on a very small scale. One brachiopod species (Coptothyrus adamsi) may be a measure of environmental conditions around an oil terminal being built in Russia on the shore of the Sea of Japan. The word "brachiopod" is formed from the Ancient Greek words brachion ("arm") and podos ("foot"). They are often known as "lamp shells", since the curved shells of the class Terebratulida resemble pottery oil-lamps.
Anatomy
Shell structure and function
Modern brachiopods range from long, and most species are about . Magellania venosa is the largest extant species. The largest brachiopods known—Gigantoproductus and Titanaria, reaching in width—occurred in the upper part of the Lower Carboniferous. Brachiopods have two valves (shell sections), which cover the dorsal (top) and ventral (bottom) surface of the animal, unlike bivalve molluscs whose shells cover the lateral surfaces (sides). The valves are unequal in size and structure, with each having its own symmetrical form rather than the two being mirror images of each other. The formation of brachiopod shells during ontogeny builds on a set of conserved genes, including homeobox genes, that are also used to form the shells of molluscs.
The brachial valve is usually smaller and bears brachia ("arms") on its inner surface. These brachia are the origin of the phylum's name, and support the lophophore, used for feeding and respiration. The pedicle valve is usually larger, and near the hinge it has an opening for the stalk-like pedicle through which most brachiopods attach themselves to the substrate. (R. C. Moore, 1952) The brachial and pedicle valves are often called the dorsal and ventral valves, respectively, but some paleontologists regard the terms "dorsal" and "ventral" as irrelevant since they believe that the "ventral" valve was formed by a folding of the upper surface under the body. The ventral ("lower") valve actually lies above the dorsal ("upper") valve when most brachiopods are oriented in life position. In many living articulate brachiopod species, both valves are convex, the surfaces often bearing growth lines and/or other ornamentation. However, inarticulate lingulids, which burrow into the seabed, have valves that are smoother, flatter and of similar size and shape. (R. C. Moore, 1952)
Articulate ("jointed") brachiopods have a tooth and socket arrangement by which the pedicle and brachial valves hinge, locking the valves against lateral displacement. Inarticulate brachiopods have no matching teeth and sockets; their valves are held together only by muscles. (R. C. Moore, 1952)
All brachiopods have adductor muscles that are set on the inside of the pedicle valve and which close the valves by pulling on the part of the brachial valve ahead of the hinge. These muscles have both "quick" fibers that close the valves in emergencies and "catch" fibers that are slower but can keep the valves closed for long periods. Articulate brachiopods open the valves by means of abductor muscles, also known as diductors, which lie further to the rear and pull on the part of the brachial valve behind the hinge. Inarticulate brachiopods use a different opening mechanism, in which muscles reduce the length of the coelom (main body cavity) and make it bulge outwards, pushing the valves apart. Both classes open the valves to an angle of about 10 degrees. The more complex set of muscles employed by inarticulate brachiopods can also operate the valves as scissors, a mechanism that lingulids use to burrow.
Each valve consists of three layers, an outer periostracum made of organic compounds and two biomineralized layers. Articulate brachiopods have an outermost periostracum made of proteins, a "primary layer" of calcite (a form of calcium carbonate) under that, and innermost a mixture of proteins and calcite. Inarticulate brachiopod shells have a similar sequence of layers, but their composition is different from that of articulated brachiopods and also varies among the classes of inarticulate brachiopods. The Terebratulida are an example of brachiopods with a punctate shell structure; the mineralized layers are perforated by tiny open canals of living tissue, extensions of the mantle called caeca, which almost reach the outside of the primary layer. These shells can contain half of the animal's living tissue. Impunctate shells are solid without any tissue inside them. Pseudopunctate shells have tubercles formed from deformations unfurling along calcite rods. They are only known from fossil forms, and were originally mistaken for calcified punctate structures.
Lingulids and discinids, which have pedicles, have a matrix of glycosaminoglycans (long, unbranched polysaccharides), in which other materials are embedded: chitin in the periostracum; apatite containing calcium phosphate in the primary biomineralized layer; and a complex mixture in the innermost layer, containing collagen and other proteins, chitinophosphate and apatite. Craniids, which have no pedicle and cement themselves directly to hard surfaces, have a periostracum of chitin and mineralized layers of calcite. Shell growth can be described as holoperipheral, mixoperipheral, or hemiperipheral. In holoperipheral growth, distinctive of craniids, new material is added at an equal rate all around the margin. In mixoperipheral growth, found in many living and extinct articulates, new material is added to the posterior region of the shell with an anterior trend, growing towards the other shell. Hemiperipheral growth, found in lingulids, is similar to mixoperipheral growth but occurs in mostly a flat plate with the shell growing forwards and outwards.
Mantle
Brachiopods, as with molluscs, have an epithelial mantle which secretes and lines the shell, and encloses the internal organs. The brachiopod body occupies only about one-third of the internal space inside the shell, nearest the hinge. The rest of the space is lined with the mantle lobes, extensions that enclose a water-filled space in which sits the lophophore. The coelom (body cavity) extends into each lobe as a network of canals, which carry nutrients to the edges of the mantle.
Relatively new cells in a groove on the edges of the mantle secrete material that extends the periostracum. These cells are gradually displaced to the underside of the mantle by more recent cells in the groove, and switch to secreting the mineralized material of the shell valves. In other words, on the edge of the valve the periostracum is extended first, and then reinforced by extension of the mineralized layers under the periostracum. In most species the edge of the mantle also bears movable bristles, often called chaetae or setae, that may help defend the animals and may act as sensors. In some brachiopods groups of chaetae help to channel the flow of water into and out of the mantle cavity.
In most brachiopods, diverticula (hollow extensions) of the mantle penetrate through the mineralized layers of the valves into the periostraca. The function of these diverticula is uncertain and it is suggested that they may be storage chambers for chemicals such as glycogen, may secrete repellents to deter organisms that stick to the shell or may help in respiration. Experiments show that a brachiopod's oxygen consumption drops if petroleum jelly is smeared on the shell, clogging the diverticula.
Lophophore
Like bryozoans and phoronids, brachiopods have a lophophore, a crown of tentacles whose cilia (fine hairs) create a water current that enables them to filter food particles out of the water. However a bryozoan or phoronid lophophore is a ring of tentacles mounted on a single, retracted stalk, while the basic form of the brachiopod lophophore is U-shaped, forming the brachia ("arms") from which the phylum gets its name. Brachiopod lophophores are non-retractable and occupy up to two-thirds of the internal space, in the frontmost area where the valves gape when opened. To provide enough filtering capacity in this restricted space, lophophores of larger brachiopods are folded in moderately to very complex shapes—loops and coils are common, and some species' lophophores contort into a shape resembling a hand with the fingers splayed. In all species the lophophore is supported by cartilage and by a hydrostatic skeleton (in other words, by the pressure of its internal fluid), and the fluid extends into the tentacles. Some articulate brachiopods also have a brachidium, a calcareous support for the lophophore attached to the inside of the brachial valve, which have led to an extremely reduced lophophoral muscles and the reduction of some brachial nerves.
The tentacles bear cilia (fine mobile hairs) on their edges and along the center. The beating of the outer cilia drives a water current from the tips of the tentacles to their bases, where it exits. Food particles that collide with the tentacles are trapped by mucus, and the cilia down the middle drive this mixture to the base of the tentacles. A brachial groove runs round the bases of the tentacles, and its own cilia pass food along the groove towards the mouth. The method used by brachiopods is known as "upstream collecting", as food particles are captured as they enter the field of cilia that creates the feeding current. This method is used by the related phoronids and bryozoans, and also by pterobranchs. Entoprocts use a similar-looking crown of tentacles, but it is solid and the flow runs from bases to tips, forming a "downstream collecting" system that catches food particles as they are about to exit.
Pedicle and other attachments
Most modern species attach to hard surfaces by means of a cylindrical pedicle ("stalk"), an extension of the body wall. This has a chitinous cuticle (non-cellular "skin") and protrudes through an opening in the hinge. However, some genera have no pedicle, such as the inarticulate Crania and the articulate Lacazella; they cement the rear of the "pedicle" (ventral) valve to a surface so that the front is slightly inclined up away from the surface. In these brachiopods, the ventral valve lacks a pedicle opening. In a few articulate genera such as Neothyris and Anakinetica, the pedicles wither as the adults grow and finally lie loosely on the surface. In these genera the shells are thickened and shaped so that the opening of the gaping valves is kept free of the sediment.
Pedicles of inarticulate species are extensions of the main coelom, which houses the internal organs. A layer of longitudinal muscles lines the epidermis of the pedicle. Members of the order Lingulida have long pedicles, which they use to burrow into soft substrates, to raise the shell to the opening of the burrow to feed, and to retract the shell when disturbed. A lingulid moves its body up and down the top two-thirds of the burrow, while the remaining third is occupied only by the pedicle, with a bulb on the end that builds a "concrete" anchor. However, the pedicles of the order Discinida are short and attach to hard surfaces.
The pedicle of articulate brachiopods has no coelom, and its homology is unclear. It is constructed from a different part of the larval body, and has a compact core composed of connective tissue. Muscles at the rear of the body can straighten, bend or even rotate the pedicle. The far end of the pedicle generally has rootlike extensions or short papillae ("bumps"), which attach to hard surfaces. However, articulate brachiopods of the genus Chlidonophora use a branched pedicle to anchor in sediment. The pedicle emerges from the pedicle valve, either through a notch in the hinge or, in species where the pedicle valve is longer than the brachial, from a hole where the pedicle valve doubles back to touch the brachial valve. Some species stand with the front end upwards, while others lie horizontal with the pedicle valve uppermost.
Some early brachiopods—for example strophomenates, kutorginates and obolellates—do not attach using their pedicle, but with an entirely different structure known as the "pedicle sheath", which has no relationship to the pedicle. This structure arises from the umbo of the pedicle valve, at the centre of the earliest (metamorphic) shell at the location of the protegulum. It is sometimes associated with a fringing plate, the colleplax.
Biology
Feeding and excretion
The water flow enters the lophophore from the sides of the open valves and exits at the front of the animal. In lingulids the entrance and exit channels are formed by groups of chaetae that function as funnels. In other brachiopods the entry and exit channels are organized by the shape of the lophophore. The lophophore captures food particles, especially phytoplankton (tiny photosynthetic organisms), and deliver them to the mouth via the brachial grooves along the bases of the tentacles. The mouth is a tiny slit at the base of the lophophore. Food passes through the mouth, muscular pharynx ("throat") and oesophagus ("gullet"), all of which are lined with cilia and cells that secrete mucus and digestive enzymes. The stomach wall has branched ceca ("pouches") where food is digested, mainly within the cells.
Nutrients are transported throughout the coelom, including the mantle lobes, by cilia. The wastes produced by metabolism are broken into ammonia, which is eliminated by diffusion through the mantle and lophophore. Brachiopods have metanephridia, used by many phyla to excrete ammonia and other dissolved wastes. However, brachiopods have no sign of the podocytes, which perform the first phase of excretion in this process, and brachiopod metanephridia appear to be used only to emit sperm and ova.
The majority of food consumed by brachiopods is digestible, with very little solid waste produced. The cilia of the lophophore can change direction to eject isolated particles of indigestible matter. If the animal encounters larger lumps of undesired matter, the cilia lining the entry channels pause and the tentacles in contact with the lumps move apart to form large gaps and then slowly use their cilia to dump the lumps onto the lining of the mantle. This has its own cilia, which wash the lumps out through the opening between the valves. If the lophophore is clogged, the adductors snap the valves sharply, which creates a "sneeze" that clears the obstructions. In some inarticulate brachiopods the digestive tract is U-shaped and ends with an anus that eliminates solids from the front of the body wall. Other inarticulate brachiopods and all articulate brachiopods have a curved gut that ends blindly, with no anus. These animals bundle solid waste with mucus and periodically "sneeze" it out, using sharp contractions of the gut muscles.
Circulation and respiration
The lophophore and mantle are the only surfaces that absorb oxygen and eliminate carbon dioxide. Oxygen seems to be distributed by the fluid of the coelom, which is circulated through the mantle and driven either by contractions of the lining of the coelom or by beating of its cilia. In some species oxygen is partly carried by the respiratory pigment hemerythrin, which is transported in coelomocyte cells. The maximum oxygen consumption of brachiopods is low, and their minimum requirement is not measurable.
Brachiopods also have colorless blood, circulated by a muscular heart lying in the dorsal part of the body above the stomach. The blood passes through vessels that extend to the front and back of the body, and branch to organs including the lophophore at the front and the gut, muscles, gonads and nephridia at the rear. The blood circulation seems not to be completely closed, and the coelomic fluid and blood must mix to a degree. The main function of the blood may be to deliver nutrients.
Nervous system and senses
The "brain" of adult articulates consists of two ganglia, one above and the other below the oesophagus. Adult inarticulates have only the lower ganglion. From the ganglia and the commissures where they join, nerves run to the lophophore, the mantle lobes and the muscles that operate the valves. The edge of the mantle has probably the greatest concentration of sensors. Although not directly connected to sensory neurons, the mantle's chaetae probably send tactile signals to receptors in the epidermis of the mantle. Many brachiopods close their valves if shadows appear above them, but the cells responsible for this are unknown. Some brachiopods have statocysts, which detect changes in the animals' position.
Reproduction and life cycle
Lifespans range from 3 to over 30 years. Adults of most species are of one sex throughout their lives. The gonads are masses of developing gametes (ova or sperm), and most species have four gonads, two in each valve. Those of articulates lie in the channels of the mantle lobes, while those of inarticulates lie near the gut. Ripe gametes float into the main coelom and then exit into the mantle cavity via the metanephridia, which open on either side of the mouth. Most species release both ova and sperm into the water, but females of some species keep the embryos in brood chambers until the larvae hatch.
The cell division in the embryo is radial (cells form in stacks of rings directly above each other), holoblastic (cells are separate, although adjoining) and regulative (the type of tissue into which a cell develops is controlled by interactions between adjacent cells, rather than rigidly within each cell). While some animals develop the mouth and anus by deepening the blastopore, a "dent" in the surface of the early embryo, the blastopore of brachiopods closes up, and their mouth and anus develop from new openings.
The larvae of lingulids (Lingulida and Discinida) are planktotrophic (feeding), and swim as plankton for months resembling miniature adults, with valves, mantle lobes, a pedicle that coils in the mantle cavity, and a small lophophore, which is used for both feeding and swimming. The larvae of craniids have no pedicle or shell. As the shell becomes heavier, the juvenile sinks to the bottom and becomes a sessile adult. The larvae of articulate species (Craniiformea and Rhynchonelliformea) are lecithotrophic (non-feeding) and live only on yolk, and remain among the plankton for only a few days. The Rhynchonelliformea larvae has three larval lobes, unlike the Craniiformea which only have two larval lobes. This type of larva has a ciliated frontmost lobe that becomes the body and lophophore, a rear lobe that becomes the pedicle, and a mantle like a skirt, with the hem towards the rear. On metamorphosing into an adult, the pedicle attaches to a surface, the front lobe develops the lophophore and other organs, and the mantle rolls up over the front lobe and starts to secrete the shell. In cold seas, brachiopod growth is seasonal and the animals often lose weight in winter. These variations in growth often form growth lines in the shells. Members of some genera have survived for a year in aquaria without food.
Taxonomy
Taxonomic history
Brachiopod fossils show great diversity in the morphology of the shells and lophophore, while the modern genera show less diversity but provide soft-bodied characteristics. Both fossils and extant species have limitations that make it difficult to produce a comprehensive classification of brachiopods based on morphology. The phylum also has experienced significant convergent evolution and reversals (in which a more recent group seems to have lost a characteristic that is seen in an intermediate group, reverting to a characteristic last seen in an older group). Hence some brachiopod taxonomists believe it is premature to define higher levels of classification such as order, and recommend instead a bottom-up approach that identifies genera and then groups these into intermediate groups.
However, other taxonomists believe that some patterns of characteristics are sufficiently stable to make higher-level classifications worthwhile, although there are different views about what the higher-level classifications should be.
The "traditional" classification was defined in 1869; two further approaches were established in the 1990s:
In the "traditional" classification, brachiopods are divided into the Articulata and Inarticulata. The Articulata have toothed hinges between the valves, while the hinges of the Inarticulata are held together only by muscles.
A classification devised in the 1990s, based on the materials of which the shells are based, united the Craniida and the "articulate" brachiopods in the Calciata, which have calcite shells. The Lingulida and Discinida, combined in the Lingulata, have shells made of chitin and calcium phosphate.
A three-part scheme, also from the 1990s, places the Craniida in a separate group of its own, the Craniiformea. The Lingulida and Discinida are grouped as Linguliformea, and the Rhynchonellida and Terebratulida as Rhynchonelliformea.
About 330 living species are recognized, grouped into over 100 genera. The great majority of modern brachiopods are rhynchonelliforms (Articulata).
Modern classification
Genetic analysis performed since the 1990s has extended the understanding of the relationship between different organisms. It is now clear the brachiopods do not belong to the Deuterostomia (such as echinoderms and chordates) as was hypothesized earlier, but should be included in the broad group Protostomia, in a subgroup now called Lophotrochozoa. Although their adult morphology seems rather different, the nucleotide sequence of the 18S rRNA indicates that the phoronids (horseshoe worms) are the closest relatives of the inarticulate brachiopods, more so than articulate brachiopods. For now, the weight of evidence is inconclusive as to the exact relations within the inarticulates. Consequently, it has been suggested to include horseshoe worms in the Brachiopoda as a class named Phoronata (B.L.Cohen & Weydmann) in addition to the Craniata and Lingulata, within the subphylum Linguliformea. The other subphylum, Rhynchonelliformea contains only one extant class, which is subdivided into the extant orders Rhynchonellida, Terebratulida and Thecideida.
Orders
This shows the taxonomy of brachiopods down to the order level, including extinct groups, which make up the majority of species. Extinct groups are indicated with a (†) symbol:
Class †Hyolitha
Order †Hyolithida
Subphylum Linguliformea
Class Lingulata
Order Lingulida
Order †Acrotretida
Order †Siphonotretida
Class †Paterinata
Order †Paterinida
Subphylum Craniiformea
Class Craniata
Order Craniida
Order †Craniopsida
Order †Trimerellida
Subphylum Rhynchonelliformea
Class †Chileata
Order †Chileida
Order †Dictyonellida
Class †Obolellata
Order †Obolellida
Order †Naukatida
Class †Kutorginata
Order †Kutorginida
Class †Strophomenata
Order †Billingsellida
Order †Strophomenida
Order †Productida
Order †Orthotetida
Class Rhynchonellata
Order Rhynchonellida
Order Terebratulida
Order Thecideida
Order †Protorthida
Order †Orthida
Order †Pentamerida
Order †Atrypida
Order †Athyridida
Order †Spiriferida
Order †Spiriferinida
Ecology
Distribution and habitat
Brachiopods are an entirely marine phylum, with no known freshwater species. Most species avoid locations with strong currents or waves, and typical sites include rocky overhangs, crevices and caves, steep slopes of continental shelves, and in deep ocean floors. However, some articulate species attach to kelp or in exceptionally sheltered sites in intertidal zones. The smallest living brachiopod, Gwynia, is only about long, and lives in between gravel grains. Rhynchonelliforms, whose larvae consume only their yolks and settle and develop quickly, are often endemic to an area and form dense populations that can reach thousands per meter. Young adults often attach to the shells of more mature ones. On the other hand, inarticulate brachiopods, whose larva swim for up to a month before settling, have wide ranges. Members of the discinoid genus Pelagodiscus have a cosmopolitan distribution.
Interactions with other organisms
Brachiopods have a low metabolic rate, between one third and one tenth of that of bivalves. While brachiopods were abundant in warm, shallow seas during the Cretaceous period, most of their former niches are now occupied by bivalves, and most now live in cold and low-light conditions.
Brachiopod shells occasionally show evidence of damage by predators, and sometimes of subsequent repair. Fish and crustaceans seem to find brachiopod flesh distasteful. The fossil record shows that drilling predators like gastropods attacked molluscs and echinoids 10 to 20 times more often than they did brachiopods, suggesting that such predators attacked brachiopods by mistake or when other prey was scarce. In waters where food is scarce, the snail Capulus ungaricus steals food from bivalves, snails, tube worms, and brachiopods.
Among brachiopods only the lingulids have been fished commercially, and only on a very small scale. It is mostly the fleshy pedicle that is eaten. Brachiopods seldom settle on artificial surfaces, probably because they are vulnerable to pollution. This may make the population of Coptothyrus adamsi useful as a measure of environmental conditions around an oil terminal being built in Russia on the shore of the Sea of Japan.
Brachiopods are the state fossil of the U.S. state of Kentucky.
Evolutionary history
Fossil record
Over 12,000 fossil species are recognized, grouped into over 5,000 genera. While the largest modern brachiopods are long, a few fossils measure up to wide. The earliest confirmed brachiopods have been found in the early Cambrian, inarticulate forms appearing first, followed soon after by articulate forms. Three unmineralized species have also been found in the Cambrian, and apparently represent two distinct groups that evolved from mineralized ancestors. The inarticulate Lingula is often called a "living fossil", as very similar genera have been found all the way back to the Ordovician. On the other hand, articulate brachiopods have produced major diversifications, and suffered severe mass extinctions—but the articulate Rhynchonellida and Terebratulida, the most diverse present-day groups, appeared at the start of the Ordovician and Carboniferous, respectively.
Since 1991 Claus Nielsen has proposed a hypothesis about the development of brachiopods, adapted in 2003 by Cohen and colleagues as a hypothesis about the earliest evolution of brachiopods. This "brachiopod fold" hypothesis suggests that brachiopods evolved from an ancestor similar to Halkieria, a slug-like animal with "chain mail" on its back and a shell at the front and rear end. The hypothesis proposes that the first brachiopod converted its shells into a pair of valves by folding the rear part of its body under its front.
However, fossils from 2007 onwards have supported a new interpretation of the Early-Cambrian tommotiids, and a new hypothesis that brachiopods evolved from tommotiids. The "armor mail" of tommotiids was well-known but not in an assembled form, and it was generally assumed that tommotiids were slug-like animals similar to Halkieria, except that tommotiids' armor was made of organophosphatic compounds while that of Halkieria was made of calcite. However, fossils of a new tommotiid, Eccentrotheca, showed an assembled mail coat that formed a tube, which would indicate a sessile animal rather than a creeping slug-like one. Eccentrotheca'''s organophosphatic tube resembled that of phoronids, sessile animals that feed by lophophores and are regarded either very close relatives or a sub-group of brachiopods. Paterimitra, another mostly assembled fossil found in 2008 and described in 2009, had two symmetrical plates at the bottom, like brachiopod valves but not fully enclosing the animal's body.
At their peak in the Paleozoic, the brachiopods were among the most abundant filter-feeders and reef-builders, and occupied other ecological niches, including swimming in the jet-propulsion style of scallops. However, after the Permian–Triassic extinction event, informally known as the "Great Dying", brachiopods recovered only a third of their former diversity. It was often thought that brachiopods were actually declining in diversity, and that in some way bivalves out-competed them. However, in 1980, Gould and Calloway produced a statistical analysis that concluded that both brachiopods and bivalves increased all the way from the Paleozoic to modern times, but bivalves increased faster; the Permian–Triassic extinction was moderately severe for bivalves but devastating for brachiopods, so that brachiopods for the first time were less diverse than bivalves and their diversity after the Permian increased from a very low base; there is no evidence that bivalves out-competed brachiopods, and short-term increases or decreases for both groups appeared synchronously. In 2007 Knoll and Bambach concluded that brachiopods were one of several groups that were most vulnerable to the Permian–Triassic extinction, as all had calcareous hard parts (made of calcium carbonate) and had low metabolic rates and weak respiratory systems.
Brachiopod fossils have been useful indicators of climate changes during the Paleozoic era. When global temperatures were low, as in much of the Ordovician, the large difference in temperature between equator and poles created different collections of fossils at different latitudes. On the other hand, warmer periods, such much of the Silurian, created smaller difference in temperatures, and all seas at the low to middle latitudes were colonized by the same few brachiopod species.
Evolutionary family tree
Deuterostomes or protostomes
From about the 1940s to the 1990s, family trees based on embryological and morphological features placed brachiopods among or as a sister group to the deuterostomes. a super-phylum that includes chordates and echinoderms. Closer examination has found difficulties in the grounds on which brachiopods were affiliated with deuterostomes:
Radial cleavage in the earliest divisions of the egg appears to be the original condition for the ancestral bilaterians, in the earliest Ecdysozoa and possibly in the earliest Eutrochozoa, a major sub-group of the Lophotrochozoa. Hence radial cleavage does not imply that brachiopods are affiliated with deuterostomes.
The traditional view is that the coelom(s) in deuterostomes and protostomes form by different process, called enterocoely and schizocoely, respectively. However, research since the early 1990s has found significant exceptions. Both types of coelom construction appear among brachiopods, and therefore do not imply that brachiopods are deuterostomes.
The terms "deuterostomes" and "protostomes" originally defined distinct ways of forming the mouth from the blastopore, a depression that appears in an early stage of the embryo. However, some "protostomes" form the mouth using a process more like that typical of deuterostomes. Hence forming the mouth via a deuterostome-like process does not imply that brachiopods are affiliated with deuterostomes.
Nielsen views the brachiopods and closely related phoronids as affiliated with the deuterostome pterobranchs because their lophophores are driven by one cilium per cell, while those of bryozoans, which he regards as protostomes, have multiple cilia per cell. However, pterobranchs are hemichordates and probably closely related to echinoderms, and there is no evidence that the latest common ancestor of pterobranchs and other hemichordates or the latest common ancestor of hemichordates and echinoderms was sessile and fed by means of tentacles.
From 1988 onwards analyses based on molecular phylogeny, which compares biochemical features such as similarities in DNA, have placed brachiopods among the Lophotrochozoa, a protostome super-phylum that includes molluscs, annelids and flatworms but excludes the other protostome super-phylum Ecdysozoa, whose members include arthropods. This conclusion is unanimous among molecular phylogeny studies that use a wide selection of genes: rDNA, Hox genes, mitochondrial protein genes, single nuclear protein genes and sets of nuclear protein genes.
Some combined studies in 2000 and 2001, using both molecular and morphological data, support brachiopods as Lophotrochozoa, while others in 1998 and 2004 concluded that brachiopods were deuterostomes.
Relationship with other lophotrochozoans
The phoronids feed with a lophophore, burrow or encrust on surfaces, and build three-layered tubes made of polysaccharide, possibly chitin, mixed with particles with seabed material. Traditionally they have been regarded as a separate phylum, but increasingly detailed molecular phylogeny studies between 1997 and 2000 have concluded that phoronids are a sub-group of brachiopods. However, an analysis in 2005 concluded that phoronids are a sub-group of bryozoans.
While all molecular phylogeny studies and half the combined studies until 2008 conclude that brachiopods are lophotrochozoans, they could not identify which lophotrochozoan phylum were the closest relatives of brachiopods—except phoronids, which are a sub-group of brachiopods. However, in 2008 two analyses found that brachiopods' closest lophotrochozoan relatives were nemertines. The authors found this surprising, since nemertines have spiral cleavage in the early stages of cell division and form a trochophore larva, while brachiopods have radial cleavage and a larva that shows no sign of having evolved from a trochophore. Another study in 2008 also concluded that brachiopods are closely related to nemertines, casting doubt on the idea that brachiopods are part of a clade Lophophorata of lophophore-feeding animals within the lophotrochozoans.
Gallery
| Biology and health sciences | Lophotrochozoa | Animals |
19827803 | https://en.wikipedia.org/wiki/Nematode | Nematode | The nematodes ( or ; ; ), roundworms or eelworms constitute the phylum Nematoda. Species in the phylum inhabit a broad range of environments. Most species are free-living, feeding on microorganisms, but many are parasitic. Parasitic worms (helminths) are the cause of soil-transmitted helminthiases.
They are classified along with arthropods, tardigrades and other moulting animals in the clade Ecdysozoa. Unlike the flatworms, nematodes have a tubular digestive system, with openings at both ends. Like tardigrades, they have a reduced number of Hox genes, but their sister phylum Nematomorpha has kept the ancestral protostome Hox genotype, which shows that the reduction has occurred within the nematode phylum.
Nematode species can be difficult to distinguish from one another. Consequently, estimates of the number of nematode species are uncertain. A 2013 survey of animal biodiversity suggested there are over 25,000. Estimates of the total number of extant species are subject to even greater variation. A widely referenced 1993 article estimated there might be over a million species of nematode. A subsequent publication challenged this claim, estimating the figure to be at least 40,000 species. Although the highest estimates (up to 100 million species) have since been deprecated, estimates supported by rarefaction curves, together with the use of DNA barcoding and the increasing acknowledgment of widespread cryptic species among nematodes, have placed the figure closer to one million species.
Nematodes have successfully adapted to nearly every ecosystem: from marine (salt) to fresh water, soils, from the polar regions to the tropics, as well as the highest to the lowest of elevations. They are ubiquitous in freshwater, marine, and terrestrial environments, where they often outnumber other animals in both individual and species counts, and are found in locations as diverse as mountains, deserts, and oceanic trenches. They are found in every part of the Earth's lithosphere, even at great depths, below the surface of the Earth in gold mines in South Africa. They represent 90% of all animals on the ocean floor. In total, 4.4 × 1020 nematodes inhabit the Earth's topsoil, or approximately 60 billion for each human, with the highest densities observed in tundra and boreal forests. Their numerical dominance, often exceeding a million individuals per square meter and accounting for about 80% of all individual animals on Earth, their diversity of lifecycles, and their presence at various trophic levels point to an important role in many ecosystems. They play crucial roles in polar ecosystems. The roughly 2,271 genera are placed in 256 families. The many parasitic forms include pathogens in most plants and animals. A third of the genera occur as parasites of vertebrates; about 35 nematode species are human parasites.
Etymology
The word nematode comes from the Modern Latin compound of nema- 'thread' (from Greek nema, genitive nematos 'thread', from the stem nein 'to spin'; cf. needle) + -odes 'like, of the nature of' (cf. -oid). The addition firstly of '-oid' and then to '-ode' renders 'threadlike'.
Taxonomy and systematics
History
In 1758, Carl Linnaeus described nematodes of a few genera including Ascaris and Dracunculus, then included in the Vermes.
The name of the group Nematoda, informally called "nematodes", came from Nematoidea, originally defined by Karl Rudolphi in 1808, from Ancient Greek νῆμα (nêma, nêmatos, 'thread') and -ειδἠς (-eidēs, 'species'). It was treated as family Nematodes by Burmeister in 1837.
At its origin, the "Nematoidea" erroneously included Nematodes and Nematomorpha, attributed by Karl Theodor Ernst von Siebold in 1843. Along with Acanthocephala, Trematoda, and Cestoidea, it formed the obsolete group Entozoa, created by Rudolphi in 1808. They were classed along with Acanthocephala in the obsolete phylum Nemathelminthes by Gegenbaur in 1859.
In 1861, Karl Moriz Diesing treated the group as order Nematoda. In 1877, the taxon Nematoidea, including the family Gordiidae (horsehair worms), was promoted to the rank of phylum by Ray Lankester. The first clear distinction between the nemas and gordiids was realized by František Vejdovsky when he named the group containing the horsehair worms the order Nematomorpha in 1886.
In 1910, Grobben proposed the phylum Aschelminthes, and the nematodes were included as class Nematoda alongside the classes Rotifera, Gastrotricha, Kinorhyncha, Priapulida, and Nematomorpha.
In 1919, Nathan Cobb proposed that nematodes should be recognized alone as a phylum. He argued they should be called "nema" in English rather than "nematodes" and defined the taxon Nemates (later emended as Nemata, Latin plural of nema), listing Nematoidea sensu restricto as a synonym.
In 1932, Potts elevated the class Nematoda to the level of phylum, leaving the name the same. Although Potts' and Cobb's classifications are equivalent, both names are used, and Nematode became a popular term in zoological science.
Phylogeny
The phylogenetic relationships of the nematodes and their close relatives among the protostomes are unresolved. Traditionally, they were held to be a lineage of their own, but in the 1990s, they were proposed to form the group Ecdysozoa together with moulting animals, such as arthropods. The identity of the closest living relatives of the Nematoda has always been considered to be well resolved. Morphological and molecular phylogenetics agree with placing the roundworms as a sister taxon to the parasitic Nematomorpha; together, they make up the Nematoida. Along with the Scalidophora (formerly Cephalorhyncha), the Nematoida form the clade Cycloneuralia, but much disagreement occurs both between and among the available morphological and molecular data. The Cycloneuralia or the Introverta—depending on the validity of the former—are often ranked as a superphylum.
Systematics
Due to the lack of knowledge regarding many nematodes, their systematics is contentious. An early and influential classification was proposed by Chitwood and Chitwood—later revised by Chitwood—who divided the phylum into two classes—Aphasmidia and Phasmidia. These were later renamed Adenophorea (gland bearers) and Secernentea (secretors), respectively. The Secernentea share several characteristics, including the presence of phasmids, a pair of sensory organs located in the lateral posterior region, and this was used as the basis for this division. This scheme was adhered to in many later classifications, though the Adenophorea were not in a uniform group.
Initial studies of incomplete DNA sequences suggested the existence of five clades:
Dorylaimida
Enoplia
Spirurina
Tylenchina
Rhabditina
The Secernentea seem to be a natural group of close relatives, while the Adenophorea appear to be a paraphyletic assemblage of roundworms that retain a good number of ancestral traits. The old Enoplia do not seem to be monophyletic, either, but do contain two distinct lineages. The old group Chromadorea seems to be another paraphyletic assemblage, with the Monhysterida representing a very ancient minor group of nematodes. Among the Secernentea, the Diplogasteria may need to be united with the Rhabditia, while the Tylenchia might be paraphyletic with the Rhabditia.
The understanding of roundworm systematics and phylogeny as of 2002 is summarised below:
Phylum Nematoda
Basal order Monhysterida
Class Dorylaimida
Class Enoplea
Class Secernentea
Subclass Diplogasteria (disputed)
Subclass Rhabditia (paraphyletic?)
Subclass Spiruria
Subclass Tylenchia (disputed)
"Chromadorea" assemblage
Later work has suggested the presence of 12 clades. In 2019, a study identified one conserved signature indel (CSI) found exclusively in members of the phylum Nematoda through comparative genetic analyses. The CSI consists of a single amino acid insertion within a conserved region of a Na(+)/H(+) exchange regulatory factor protein NRFL-1 and is a molecular marker that distinguishes the phylum from other species. An analysis of the mitochondrial DNA suggests that the following groupings are valid
subclass Dorylaimia
orders Rhabditida, Trichinellida and Mermithida
suborder Rhabditina
infraorders Spiruromorpha and Oxyuridomorpha
In 2022 a new classification of the entire phylum Nematoda was presented by M. Hodda. It was based on current molecular, developmental and morphological evidence. Under this classification, the classes and subclasses are:
Class Enoplea
Subclass Enoplia
Subclass Oncholaimia
Subclass Triplonchia
Class Dorylaimida
Subclass Dorylaimia
Subclass Bathyodontia
Subclass Trichocephalia
Class Chromadorea
Subclass Chromadoria
Subclass Plectia
Fossil record
Nematode eggs from the clades Ascaridina, Spirurina, and Trichocephalida have been discovered in coprolites from the Oligocene-aged Tremembé Formation, which represented a palaeolake in present-day São Paulo with a diverse fossil assemblage of birds, fish, and arthropods that lent itself to fostering high nematode diversity. Nematodes have also been found in various lagerstätten, such as Burmese amber, the Moltrasio Formation, and the Rhynie chert, where the earliest known fossils are known from.
Anatomy
Nematodes are very small, slender worms. Most are free-living, often less than 2.5 mm long and some only about 1 mm. Many nematodes are microscopic. Some soil nematodes can reach up to 7 mm in length, and some marine species can reach up to 5 cm. Some are parasitic and can reach lengths of 50 cm or more.
The body is often ornamented with ridges, rings, bristles, or other distinctive structures.
The head is relatively distinct. Whereas the rest of the body is bilaterally symmetrical, the head is radially symmetrical, with sensory bristles and, in many cases, solid 'head-shields' radiating outwards around the mouth. The mouth has either three or six lips, which often bear a series of teeth on their inner edges. An adhesive 'caudal gland' is often found at the tip of the tail. The epidermis is either a syncytium or a single layer of cells, and is covered by a thick collagenous cuticle. The cuticle is often of a complex structure and may have two or three distinct layers. Underneath the epidermis lies a layer of longitudinal muscle cells. The relatively rigid cuticle works with the muscles to create a hydroskeleton, as nematodes lack circumferential muscles. Projections run from the inner surface of muscle cells towards the nerve cords; this is a unique arrangement in the animal kingdom, in which nerve cells normally extend fibers into the muscles rather than vice versa.
Digestive system
The oral cavity is lined with cuticles, which are often strengthened with structures, such as ridges, especially in carnivorous species, which may bear several teeth. The mouth often includes a sharp stylet, which the animal can thrust into its prey. In some species, the stylet is hollow and can be used to suck liquids from plants or animals. The oral cavity opens into a muscular, sucking pharynx, also lined with cuticle. Digestive glands are found in this region of the gut, producing enzymes that start to break down the food. In stylet-bearing species, these may even be injected into the prey.
No stomach is present, with the pharynx connecting directly to a muscleless intestine that forms the main length of the gut. This produces further enzymes and also absorbs nutrients through its single-cell-thick lining. The last portion of the intestine is lined by a cuticle, forming a rectum, which expels waste through the anus just below and in front of the tip of the tail. The movement of food through the digestive system is the result of the body movements of the worm. The intestine has valves or sphincters at either end to help control food movement through the body.
Excretory system
Nitrogenous waste is excreted in the form of ammonia through the body wall, and is not associated with any specific organs. However, the structures for excreting salt to maintain osmoregulation are typically more complex.
There is an excretory gland, also known as a ventral cell, or renette cell in all species of Adenophorea. In Secernentia there is an excretory canal system that may or may not use a gland cell.
Nervous system
At the anterior end of the animal a dense, circular nerve ring which serves as the brain surrounds the pharynx. From this ring six labial papillary nerve cords extend anteriorly, while six nerve cords; a large ventral, a smaller dorsal and two pairs of sublateral cords extend posteriorly. Each nerve lies within a cord of connective tissue lying beneath the cuticle and between the muscle cells. The ventral nerve is the largest, and has a double structure forward of the excretory pore. The dorsal nerve is responsible for motor control, while the lateral nerves are sensory, and the ventral combines both functions.
The nervous system is the only place in the body that contains cilia; these are all nonmotile and with a sensory function.
The body is covered in numerous sensory bristles and papillae that together provide a sense of touch. Behind the sensory bristles on the head lie two small pits, or 'amphids'. These are well supplied with nerve cells and are probably chemoreception organs. A few aquatic nematodes possess what appear to be pigmented eye-spots, but whether or not these are actually sensory in nature is unclear.
Reproduction
Most nematode species are dioecious, with separate male and female individuals, though some, such as Caenorhabditis elegans, are androdioecious, consisting of hermaphrodites and rare males. Both sexes possess one or two tubular gonads. In males, the sperm are produced at the end of the gonad and migrate along its length as they mature. The testis opens into a relatively wide seminal vesicle and then during intercourse into a glandular and muscular ejaculatory duct associated with the vas deferens and cloaca. In females, the ovaries each open into an oviduct (in hermaphrodites, the eggs enter a spermatheca first) and then a glandular uterus. The uteri both open into a common vulva/vagina, usually located in the middle of the morphologically ventral surface.
Reproduction is usually sexual, though hermaphrodites are capable of self-fertilization. Males are usually smaller than females or hermaphrodites (often much smaller) and often have a characteristically bent or fan-shaped tail. During copulation, one or more chitinized spicules move out of the cloaca and are inserted into the genital pore of the female. Amoeboid sperm crawl along the spicule into the female worm. Nematode sperm is thought to be the only eukaryotic cell without the globular protein G-actin.
Eggs may be embryonated or unembryonated when passed by the female, meaning their fertilized eggs may not yet be developed. A few species are known to be ovoviviparous. The eggs are protected by an outer shell, secreted by the uterus. In free-living roundworms, the eggs hatch into larvae, which appear essentially identical to the adults, except for an underdeveloped reproductive system; in parasitic roundworms, the lifecycle is often much more complicated. The structure of the eggshell is complicated and includes several layers; a detailed anatomical and terminological framework has been proposed for these layers in 2023.
Nematodes as a whole possess a wide range of modes of reproduction. Some nematodes, such as Heterorhabditis spp., undergo a process called endotokia matricida: intrauterine birth causing maternal death. Some nematodes are hermaphroditic, and keep their self-fertilized eggs inside the uterus until they hatch. The juvenile nematodes then ingest the parent nematode. This process is significantly promoted in environments with a low food supply.
The nematode model species C. elegans, C. briggsae, and Pristionchus pacificus, among other species, exhibit androdioecy,
> which is otherwise very rare among animals. The single genus Meloidogyne (root-knot nematodes) exhibits a range of reproductive modes, including sexual reproduction, facultative sexuality (in which most, but not all, generations reproduce asexually), and both meiotic and mitotic parthenogenesis.
The genus Mesorhabditis exhibits an unusual form of parthenogenesis, in which sperm-producing males copulate with females, but the sperm do not fuse with the ovum. Contact with the sperm is essential for the ovum to begin dividing, but because no fusion of the cells occurs, the male contributes no genetic material to the offspring, which are essentially clones of the female.
Aging
The nematode Caenorhabditis elegans is often used as a model organism for studying aging at the molecular level. For example, in C. elegans aging negatively impacts DNA repair, and mutants of C. elegans that are long-lived were shown to have increased DNA repair capability. These findings suggest a genetically determined correlation between DNA repair capacity and lifespan. In female C. elegans, germline processes that control DNA repair and formation of chromosomal crossovers during meiosis were shown to progressively deteriorate with age.
Free-living species
Different free-living species feed on materials as varied as bacteria, algae, fungi, small animals, fecal matter, dead organisms, and living tissues. Free-living marine nematodes are important and abundant members of the meiobenthos. They play an important role in the decomposition process, aid in recycling of nutrients in marine environments, and are sensitive to changes in the environment caused by pollution. One roundworm of note, C. elegans, lives in the soil and has found much use as a model organism. C. elegans has had its entire genome sequenced, the developmental fate of every cell determined, and every neuron mapped.
Parasitic species
Nematodes that commonly parasitise humans include ascarids (Ascaris), filarias, hookworms, pinworms (Enterobius), and whipworms (Trichuris trichiura). The species Trichinella spiralis, commonly known as the trichina worm, occurs in rats, pigs, bears, and humans, and is responsible for the disease trichinosis. Baylisascaris usually infests wild animals, but can be deadly to humans, as well. Dirofilaria immitis is known for causing heartworm disease by inhabiting the hearts, arteries, and lungs of dogs and some cats. Haemonchus contortus is one of the most abundant infectious agents in sheep around the world, causing great economic damage to sheep. In contrast, entomopathogenic nematodes parasitize insects and are mostly considered beneficial by humans, but some attack beneficial insects.
One form of nematode is entirely dependent upon fig wasps, which are the sole source of fig fertilization. They prey upon the wasps, riding them from the ripe fig of the wasp's birth to the fig flower of its death, where they kill the wasp, and their offspring await the birth of the next generation of wasps as the fig ripens.
A parasitic tetradonematid nematode discovered in 2005, Myrmeconema neotropicum, induces fruit mimicry in the tropical ant Cephalotes atratus. Infected ants develop bright red gasters (abdomens), tend to be more sluggish, and walk with their gasters in a conspicuous elevated position. These changes likely cause frugivorous birds to confuse the infected ants for berries, and eat them. Parasite eggs passed in the bird's feces are subsequently collected by foraging C. atratus and are fed to their larvae, thus completing the lifecycle of M. neotropicum.
Similarly, multiple varieties of nematodes have been found in the abdominal cavities of the primitively social sweat bee, Lasioglossum zephyrus. Inside the female body, the nematode hinders ovarian development and renders the bee less active, thus less effective in pollen collection.
Agriculture and horticulture
Depending on its species, a nematode may be beneficial or detrimental to plant health. From agricultural and horticulture perspectives, the two categories of nematodes are the predatory ones, which kill garden pests; and the pest nematodes, which attack plants, or act as vectors spreading plant viruses between crop plants. Predatory nematodes include Phasmarhabditis hermaphrodita which is a lethal parasite of gastropods such as slugs and snails. Some members of the genus Steinernema such as Steinernema carpocapsae and Steinernema riobrave are generalist parasites of webworms, cutworms, armyworms, girdlers, some weevils, wood-borers and corn earworm moths. These organisms are grown commercially as biological pest control agents which can be used as an alternative to pesticides; their use is considered very safe. Plant-parasitic nematodes include several groups causing severe crop losses, taking 10% of crops worldwide every year. The most common genera are Aphelenchoides (foliar nematodes), Ditylenchus, Globodera (potato cyst nematodes), Heterodera (soybean cyst nematodes), Longidorus, Meloidogyne (root-knot nematodes), Nacobbus, Pratylenchus (lesion nematodes), Trichodorus, and Xiphinema (dagger nematodes). Several phytoparasitic nematode species cause histological damages to roots, including the formation of visible galls (e.g. by root-knot nematodes), which are useful characters for their diagnostic in the field. Some nematode species transmit plant viruses through their feeding activity on roots. One of them is Xiphinema index, vector of grapevine fanleaf virus, an important disease of grapes, another one is Xiphinema diversicaudatum, vector of arabis mosaic virus. Other nematodes attack bark and forest trees. The most important representative of this group is Bursaphelenchus xylophilus, the pine wood nematode, present in Asia and America and recently discovered in Europe. This nematode is transmitted from tree to tree by sawyer beetles (Monochamus).
Greenhouse growers use entomopathogenic nematodes as beneficial agents to control fungus gnats. The nematodes enter the larvae of the gnats by way of their anus, mouth, and spiracles (breathing pores) and then release bacteria which kills the gnat larvae. Commonly used nematode species to control pests on greenhouse crops include Steinernema feltiae for fungus gnats and western flower thrips, Steinernema carpocapsae used to control shore flies, Steinernema kraussei for control of black vine weevils, and Heterorhabditis bacteriophora to control beetle larvae.
Rotations of plants with nematode-resistant species or varieties is one means of managing parasitic nematode infestations. For example, planting Tagetes marigolds as a cover crop just prior to planting a nematode-susceptible plant, has been shown to suppress nematodes. Another is treatment with natural antagonists such as the fungus Gliocladium roseum. Chitosan, a natural biocontrol, elicits plant defense responses to destroy parasitic cyst nematodes on roots of soybean, corn, sugar beet, potato, and tomato crops without harming beneficial nematodes in the soil. Soil steaming is an efficient method to kill nematodes before planting a crop, but indiscriminately eliminates both harmful and beneficial soil fauna.
The golden nematode Globodera rostochiensis is a particularly harmful pest that has resulted in quarantines and crop failures worldwide. It can be controlled, however. CSIRO, the scientific research body of the Australian government, found a 13- to 14-fold reduction of nematode population densities in plots having Chinese mustard Brassica juncea green manure or seed meal in the soil.
Disease in humans
A number of pathogenic intestinal nematodes cause diseases in humans, including ascariasis, trichuriasis, and hookworm disease. Anisakis species parasitise fish, and marine mammals and when consumed by humans can cause anisakiasis a gastric or gastroallergic disease. Gastrointestinal nematode infections in humans are common, with approximately 50% of the global population being affected. Developing countries are most heavily impacted, in part due to lack of access to medical care.
Trichinosis starts in the intestines but larvae can migrate to muscle. Filarial nematodes cause filariases.
Toxocariasis is a zoonotic infection caused by roundworms passed from dogs, and sometimes cats. It can give rise to different types of larva migrans such as visceral larva migrans, and ocular larva migrans.
Studies have shown that parasitic nematodes infect American eels causing damage to the eel's swim bladder, dairy animals like cattle and buffalo, and all species of sheep.
Soil ecosystems
About 90% of nematodes reside in the top 15 cm (6") of soil. Nematodes do not decompose organic matter, but, instead, are parasitic and free-living organisms that feed on living material. Nematodes can effectively regulate bacterial population and community composition—they may eat up to 5,000 bacteria per minute. Also, nematodes can play an important role in the nitrogen cycle by way of nitrogen mineralization. But plant parasitic nematodes cause billions of dollars in annual crop damage worldwide.
One group of carnivorous fungi, the nematophagous fungi, are predators of soil nematodes. They can set enticements for the nematodes in the form of lassos or adhesive structures. They can also release powerful toxins when in contact with nematodes.
Survivability
The nematode Caenorhabditis elegans an important model organism, was used as part of an ongoing research project conducted on the 2003 Space Shuttle Columbia mission STS-107, and survived the re-entry breakup. It is believed to be the first known species to survive a virtually unprotected atmospheric descent to Earth's surface. The Antarctic nematode Panagrolaimus davidi was able to withstand intracellular freezing depending on how well it had been fed. In 2023 an individual of Panagrolaimus kolymaensis was revived after 46,000 years in Siberian permafrost.
| Biology and health sciences | Ecdysozoa | null |
19828134 | https://en.wikipedia.org/wiki/Plant | Plant | Plants are the eukaryotes that form the kingdom Plantae; they are predominantly photosynthetic. This means that they obtain their energy from sunlight, using chloroplasts derived from endosymbiosis with cyanobacteria to produce sugars from carbon dioxide and water, using the green pigment chlorophyll. Exceptions are parasitic plants that have lost the genes for chlorophyll and photosynthesis, and obtain their energy from other plants or fungi. Most plants are multicellular, except for some green algae.
Historically, as in Aristotle's biology, the plant kingdom encompassed all living things that were not animals, and included algae and fungi. Definitions have narrowed since then; current definitions exclude the fungi and some of the algae. By the definition used in this article, plants form the clade Viridiplantae (green plants), which consists of the green algae and the embryophytes or land plants (hornworts, liverworts, mosses, lycophytes, ferns, conifers and other gymnosperms, and flowering plants). A definition based on genomes includes the Viridiplantae, along with the red algae and the glaucophytes, in the clade Archaeplastida.
There are about 380,000 known species of plants, of which the majority, some 260,000, produce seeds. They range in size from single cells to the tallest trees. Green plants provide a substantial proportion of the world's molecular oxygen; the sugars they create supply the energy for most of Earth's ecosystems, and other organisms, including animals, either eat plants directly or rely on organisms which do so.
Grain, fruit, and vegetables are basic human foods and have been domesticated for millennia. People use plants for many purposes, such as building materials, ornaments, writing materials, and, in great variety, for medicines. The scientific study of plants is known as botany, a branch of biology.
Definition
Taxonomic history
All living things were traditionally placed into one of two groups, plants and animals. This classification dates from Aristotle (384–322 BC), who distinguished different levels of beings in his biology, based on whether living things had a "sensitive soul" or like plants only a "vegetative soul". Theophrastus, Aristotle's student, continued his work in plant taxonomy and classification. Much later, Linnaeus (1707–1778) created the basis of the modern system of scientific classification, but retained the animal and plant kingdoms, naming the plant kingdom the Vegetabilia.
Alternative concepts
When the name Plantae or plant is applied to a specific group of organisms or taxa, it usually refers to one of four concepts. From least to most inclusive, these four groupings are:
Evolution
Diversity
There are about 382,000 accepted species of plants, of which the great majority, some 283,000, produce seeds. The table below shows some species count estimates of different green plant (Viridiplantae) divisions. About 85–90% of all plants are flowering plants. Several projects are currently attempting to collect records on all plant species in online databases, e.g. the World Flora Online.
Plants range in scale from single-celled organisms such as desmids (from across) and picozoa (less than across), to the largest trees (megaflora) such as the conifer Sequoia sempervirens (up to tall) and the angiosperm Eucalyptus regnans (up to tall).
The naming of plants is governed by the International Code of Nomenclature for algae, fungi, and plants and the International Code of Nomenclature for Cultivated Plants.
Evolutionary history
The ancestors of land plants evolved in water. An algal scum formed on the land , but it was not until the Ordovician, around , that the first land plants appeared, with a level of organisation like that of bryophytes. However, fossils of organisms with a flattened thallus in Precambrian rocks suggest that multicellular freshwater eukaryotes existed over 1000 mya.
Primitive land plants began to diversify in the late Silurian, around . Bryophytes, club mosses, and ferns then appear in the fossil record. Early plant anatomy is preserved in cellular detail in an early Devonian fossil assemblage from the Rhynie chert. These early plants were preserved by being petrified in chert formed in silica-rich volcanic hot springs.
By the end of the Devonian, most of the basic features of plants today were present, including roots, leaves and secondary wood in trees such as Archaeopteris. The Carboniferous period saw the development of forests in swampy environments dominated by clubmosses and horsetails, including some as large as trees, and the appearance of early gymnosperms, the first seed plants. The Permo-Triassic extinction event radically changed the structures of communities. This may have set the scene for the evolution of flowering plants in the Triassic (~), with an adaptive radiation in the Cretaceous so rapid that Darwin called it an "abominable mystery". Conifers diversified from the Late Triassic onwards, and became a dominant part of floras in the Jurassic.
Phylogeny
In 2019, a phylogeny based on genomes and transcriptomes from 1,153 plant species was proposed. The placing of algal groups is supported by phylogenies based on genomes from the Mesostigmatophyceae and Chlorokybophyceae that have since been sequenced. Both the "chlorophyte algae" and the "streptophyte algae" are treated as paraphyletic (vertical bars beside phylogenetic tree diagram) in this analysis, as the land plants arose from within those groups. The classification of Bryophyta is supported both by Puttick et al. 2018, and by phylogenies involving the hornwort genomes that have also since been sequenced.
Physiology
Plant cells
Plant cells have distinctive features that other eukaryotic cells (such as those of animals) lack. These include the large water-filled central vacuole, chloroplasts, and the strong flexible cell wall, which is outside the cell membrane. Chloroplasts are derived from what was once a symbiosis of a non-photosynthetic cell and photosynthetic cyanobacteria. The cell wall, made mostly of cellulose, allows plant cells to swell up with water without bursting. The vacuole allows the cell to change in size while the amount of cytoplasm stays the same.
Plant structure
Most plants are multicellular. Plant cells differentiate into multiple cell types, forming tissues such as the vascular tissue with specialized xylem and phloem of leaf veins and stems, and organs with different physiological functions such as roots to absorb water and minerals, stems for support and to transport water and synthesized molecules, leaves for photosynthesis, and flowers for reproduction.
Photosynthesis
Plants photosynthesize, manufacturing food molecules (sugars) using energy obtained from light. Plant cells contain chlorophylls inside their chloroplasts, which are green pigments that are used to capture light energy. The end-to-end chemical equation for photosynthesis is:
6CO2{} + 6H2O{} ->[\text{light}] C6H12O6{} + 6O2{}
This causes plants to release oxygen into the atmosphere. Green plants provide a substantial proportion of the world's molecular oxygen, alongside the contributions from photosynthetic algae and cyanobacteria.
Plants that have secondarily adopted a parasitic lifestyle may lose the genes involved in photosynthesis and the production of chlorophyll.
Growth and repair
Growth is determined by the interaction of a plant's genome with its physical and biotic environment. Factors of the physical or abiotic environment include temperature, water, light, carbon dioxide, and nutrients in the soil. Biotic factors that affect plant growth include crowding, grazing, beneficial symbiotic bacteria and fungi, and attacks by insects or plant diseases.
Frost and dehydration can damage or kill plants. Some plants have antifreeze proteins, heat-shock proteins and sugars in their cytoplasm that enable them to tolerate these stresses. Plants are continuously exposed to a range of physical and biotic stresses which cause DNA damage, but they can tolerate and repair much of this damage.
Reproduction
Plants reproduce to generate offspring, whether sexually, involving gametes, or asexually, involving ordinary growth. Many plants use both mechanisms.
Sexual
When reproducing sexually, plants have complex lifecycles involving alternation of generations. One generation, the sporophyte, which is diploid (with 2 sets of chromosomes), gives rise to the next generation, the gametophyte, which is haploid (with one set of chromosomes). Some plants also reproduce asexually via spores. In some non-flowering plants such as mosses, the sexual gametophyte forms most of the visible plant. In seed plants (gymnosperms and flowering plants), the sporophyte forms most of the visible plant, and the gametophyte is very small. Flowering plants reproduce sexually using flowers, which contain male and female parts: these may be within the same (hermaphrodite) flower, on different flowers on the same plant, or on different plants. The stamens create pollen, which produces male gametes that enter the ovule to fertilize the egg cell of the female gametophyte. Fertilization takes place within the carpels or ovaries, which develop into fruits that contain seeds. Fruits may be dispersed whole, or they may split open and the seeds dispersed individually.
Asexual
Plants reproduce asexually by growing any of a wide variety of structures capable of growing into new plants. At the simplest, plants such as mosses or liverworts may be broken into pieces, each of which may regrow into whole plants. The propagation of flowering plants by cuttings is a similar process. Structures such as runners enable plants to grow to cover an area, forming a clone. Many plants grow food storage structures such as tubers or bulbs which may each develop into a new plant.
Some non-flowering plants, such as many liverworts, mosses and some clubmosses, along with a few flowering plants, grow small clumps of cells called gemmae which can detach and grow.
Disease resistance
Plants use pattern-recognition receptors to recognize pathogens such as bacteria that cause plant diseases. This recognition triggers a protective response. The first such plant receptors were identified in rice and in Arabidopsis thaliana.
Genomics
Plants have some of the largest genomes of all organisms. The largest plant genome (in terms of gene number) is that of wheat (Triticum aestivum), predicted to encode ≈94,000 genes and thus almost 5 times as many as the human genome. The first plant genome sequenced was that of Arabidopsis thaliana which encodes about 25,500 genes. In terms of sheer DNA sequence, the smallest published genome is that of the carnivorous bladderwort (Utricularia gibba) at 82 Mb (although it still encodes 28,500 genes) while the largest, from the Norway spruce (Picea abies), extends over 19.6 Gb (encoding about 28,300 genes).
Ecology
Distribution
Plants are distributed almost worldwide. While they inhabit many biomes which can be divided into a multitude of ecoregions, only the hardy plants of the Antarctic flora, consisting of algae, mosses, liverworts, lichens, and just two flowering plants, have adapted to the prevailing conditions on that southern continent.
Plants are often the dominant physical and structural component of the habitats where they occur. Many of the Earth's biomes are named for the type of vegetation because plants are the dominant organisms in those biomes, such as grassland, savanna, and tropical rainforest.
Primary producers
The photosynthesis conducted by land plants and algae is the ultimate source of energy and organic material in nearly all ecosystems. Photosynthesis, at first by cyanobacteria and later by photosynthetic eukaryotes, radically changed the composition of the early Earth's anoxic atmosphere, which as a result is now 21% oxygen. Animals and most other organisms are aerobic, relying on oxygen; those that do not are confined to relatively rare anaerobic environments. Plants are the primary producers in most terrestrial ecosystems and form the basis of the food web in those ecosystems. Plants form about 80% of the world biomass at about of carbon.
Ecological relationships
Numerous animals have coevolved with plants; flowering plants have evolved pollination syndromes, suites of flower traits that favour their reproduction. Many, including insect and bird partners, are pollinators, visiting flowers and accidentally transferring pollen in exchange for food in the form of pollen or nectar.
Many animals disperse seeds that are adapted for such dispersal. Various mechanisms of dispersal have evolved. Some fruits offer nutritious outer layers attractive to animals, while the seeds are adapted to survive the passage through the animal's gut; others have hooks that enable them to attach to a mammal's fur.
Myrmecophytes are plants that have coevolved with ants. The plant provides a home, and sometimes food, for the ants. In exchange, the ants defend the plant from herbivores and sometimes competing plants. Ant wastes serve as organic fertilizer.
The majority of plant species have fungi associated with their root systems in a mutualistic symbiosis known as mycorrhiza. The fungi help the plants gain water and mineral nutrients from the soil, while the plant gives the fungi carbohydrates manufactured in photosynthesis.
Some plants serve as homes for endophytic fungi that protect the plant from herbivores by producing toxins. The fungal endophyte Neotyphodium coenophialum in tall fescue grass has pest status in the American cattle industry.
Many legumes have Rhizobium nitrogen-fixing bacteria in nodules of their roots, which fix nitrogen from the air for the plant to use; in return, the plants supply sugars to the bacteria. Nitrogen fixed in this way can become available to other plants, and is important in agriculture; for example, farmers may grow a crop rotation of a legume such as beans, followed by a cereal such as wheat, to provide cash crops with a reduced input of nitrogen fertilizer.
Some 1% of plants are parasitic. They range from the semi-parasitic mistletoe that merely takes some nutrients from its host, but still has photosynthetic leaves, to the fully-parasitic broomrape and toothwort that acquire all their nutrients through connections to the roots of other plants, and so have no chlorophyll. Full parasites can be extremely harmful to their plant hosts.
Plants that grow on other plants, usually trees, without parasitizing them, are called epiphytes. These may support diverse arboreal ecosystems. Some may indirectly harm their host plant, such as by intercepting light. Hemiepiphytes like the strangler fig begin as epiphytes, but eventually set their own roots and overpower and kill their host. Many orchids, bromeliads, ferns, and mosses grow as epiphytes. Among the epiphytes, the bromeliads accumulate water in their leaf axils; these water-filled cavities can support complex aquatic food webs.
Some 630 species of plants are carnivorous, such as the Venus flytrap (Dionaea muscipula) and sundew (Drosera species). They trap small animals and digest them to obtain mineral nutrients, especially nitrogen and phosphorus.
Competition
Competition for shared resources reduces a plant's growth. Shared resources include sunlight, water and nutrients. Light is a critical resource because it is necessary for photosynthesis. Plants use their leaves to shade other plants from sunlight and grow quickly to maximize their own expose. Water too is essential for photosynthesis; roots compete to maximize water uptake from soil. Some plants have deep roots that are able to locate water stored deep underground, and others have shallower roots that are capable of extending longer distances to collect recent rainwater.
Minerals are important for plant growth and development. Common nutrients competed for amongst plants include nitrogen, phosphorus, and potassium.
Importance to humans
Food
Human cultivation of plants is the core of agriculture, which in turn has played a key role in the history of world civilizations. Humans depend on flowering plants for food, either directly or as feed in animal husbandry. More broadly, agriculture includes agronomy for arable crops, horticulture for vegetables and fruit, and forestry, including both flowering plants and conifers, for timber. About 7,000 species of plant have been used for food, though most of today's food is derived from only 30 species. The major staples include cereals such as rice and wheat, starchy roots and tubers such as cassava and potato, and legumes such as peas and beans. Vegetable oils such as olive oil and palm oil provide lipids, while fruit and vegetables contribute vitamins and minerals to the diet. Coffee, tea, and chocolate are major crops whose caffeine-containing products serve as mild stimulants. The study of plant uses by people is called economic botany or ethnobotany.
Medicines
Medicinal plants are a primary source of organic compounds, both for their medicinal and physiological effects, and for the industrial synthesis of a vast array of organic chemicals. Many hundreds of medicines, as well as narcotics, are derived from plants, both traditional medicines used in herbalism and chemical substances purified from plants or first identified in them, sometimes by ethnobotanical search, and then synthesised for use in modern medicine. Modern medicines derived from plants include aspirin, taxol, morphine, quinine, reserpine, colchicine, digitalis and vincristine. Plants used in herbalism include ginkgo, echinacea, feverfew, and Saint John's wort. The pharmacopoeia of Dioscorides, , describing some 600 medicinal plants, was written between 50 and 70 CE and remained in use in Europe and the Middle East until around 1600 CE; it was the precursor of all modern pharmacopoeias.
Nonfood products
Plants grown as industrial crops are the source of a wide range of products used in manufacturing. Nonfood products include essential oils, natural dyes, pigments, waxes, resins, tannins, alkaloids, amber and cork. Products derived from plants include soaps, shampoos, perfumes, cosmetics, paint, varnish, turpentine, rubber, latex, lubricants, linoleum, plastics, inks, and gums. Renewable fuels from plants include firewood, peat and other biofuels. The fossil fuels coal, petroleum and natural gas are derived from the remains of aquatic organisms including phytoplankton in geological time. Many of the coal fields date to the Carboniferous period of Earth's history. Terrestrial plants also form type III kerogen, a source of natural gas.
Structural resources and fibres from plants are used to construct dwellings and to manufacture clothing. Wood is used for buildings, boats, and furniture, and for smaller items such as musical instruments and sports equipment. Wood is pulped to make paper and cardboard. Cloth is often made from cotton, flax, ramie or synthetic fibres such as rayon, derived from plant cellulose. Thread used to sew cloth likewise comes in large part from cotton.
Ornamental plants
Thousands of plant species are cultivated for their beauty and to provide shade, modify temperatures, reduce wind, abate noise, provide privacy, and reduce soil erosion. Plants are the basis of a multibillion-dollar per year tourism industry, which includes travel to historic gardens, national parks, rainforests, forests with colourful autumn leaves, and festivals such as Japan's and America's cherry blossom festivals.
Plants may be grown indoors as houseplants, or in specialized buildings such as greenhouses. Plants such as Venus flytrap, sensitive plant and resurrection plant are sold as novelties. Art forms specializing in the arrangement of cut or living plant include bonsai, ikebana, and the arrangement of cut or dried flowers. Ornamental plants have sometimes changed the course of history, as in tulipomania.
In science
The traditional study of plants is the science of botany. Basic biological research has often used plants as its model organisms. In genetics, the breeding of pea plants allowed Gregor Mendel to derive the basic laws governing inheritance, and examination of chromosomes in maize allowed Barbara McClintock to demonstrate their connection to inherited traits. The plant Arabidopsis thaliana is used in laboratories as a model organism to understand how genes control the growth and development of plant structures. Tree rings provide a method of dating in archeology, and a record of past climates. The study of plant fossils, or Paleobotany, provides information about the evolutions of plants, paleogeographical reconstructions, and past climate change. Plant fossils can also help determine the age of rocks.
In mythology, religion, and culture
Plants including trees appear in mythology, religion, and literature. In multiple Indo-European, Siberian, and Native American religions, the world tree motif is depicted as a colossal tree growing on the earth, supporting the heavens, and with its roots reaching into the underworld. It may also appear as a cosmic tree or an eagle and serpent tree. Forms of the world tree include the archetypal tree of life, which is in turn connected to the Eurasian concept of the sacred tree. Another widespread ancient motif, found for example in Iran, has a tree of life flanked by a pair of confronted animals.
Flowers are often used as memorials, gifts and to mark special occasions such as births, deaths, weddings and holidays. Flower arrangements may be used to send hidden messages. Plants and especially flowers form the subjects of many paintings.
Negative effects
Weeds are commercially or aesthetically undesirable plants growing in managed environments such as in agriculture and gardens. People have spread many plants beyond their native ranges; some of these plants have become invasive, damaging existing ecosystems by displacing native species, and sometimes becoming serious weeds of cultivation.
Some plants that produce windblown pollen, including grasses, invoke allergic reactions in people who suffer from hay fever. Many plants produce toxins to protect themselves from herbivores. Major classes of plant toxins include alkaloids, terpenoids, and phenolics. These can be harmful to humans and livestock by ingestion or, as with poison ivy, by contact. Some plants have negative effects on other plants, preventing seedling growth or the growth of nearby plants by releasing allopathic chemicals.
| Biology and health sciences | Science and medicine | null |
6354378 | https://en.wikipedia.org/wiki/Drug%20delivery | Drug delivery | Drug delivery refers to approaches, formulations, manufacturing techniques, storage systems, and technologies involved in transporting a pharmaceutical compound to its target site to achieve a desired therapeutic effect. Principles related to drug preparation, route of administration, site-specific targeting, metabolism, and toxicity are used to optimize efficacy and safety, and to improve patient convenience and compliance. Drug delivery is aimed at altering a drug's pharmacokinetics and specificity by formulating it with different excipients, drug carriers, and medical devices. There is additional emphasis on increasing the bioavailability and duration of action of a drug to improve therapeutic outcomes. Some research has also been focused on improving safety for the person administering the medication. For example, several types of microneedle patches have been developed for administering vaccines and other medications to reduce the risk of needlestick injury.
Drug delivery is a concept heavily integrated with dosage form and route of administration, the latter sometimes being considered part of the definition. While route of administration is often used interchangeably with drug delivery, the two are separate concepts. Route of administration refers to the path a drug takes to enter the body, whereas drug delivery also encompasses the engineering of delivery systems and can include different dosage forms and devices used to deliver a drug through the same route. Common routes of administration include oral, parenteral (injected), sublingual, topical, transdermal, nasal, ocular, rectal, and vaginal, however, drug delivery is not limited to these routes and there may be several ways to deliver medications through other routes.
Since the approval of the first controlled-release formulation in the 1950s, research into new delivery systems has been progressing, as opposed to new drug development which has been declining. Several factors may be contributing to this shift in focus. One of the driving factors is the high cost of developing new drugs. A 2013 review found the cost of developing a delivery system was only 10% of the cost of developing a new pharmaceutical. A more recent study found the median cost of bringing a new drug to market was $985 million in 2020, but did not look at the cost of developing drug delivery systems. Other factors that have potentially influenced the increase in drug delivery system development may include the increasing prevalence of both chronic and infectious diseases, as well as a general increased understanding of the pharmacology, pharmacokinetics, and pharmacodynamics of many drugs.
Current efforts
Current efforts in drug delivery are vast and include topics such as controlled-release formulations, targeted delivery, nanomedicine, drug carriers, 3D printing, and the delivery of biologic drugs.
The relation between nanomaterial and drug delivery
Nanotechnology is a broad field of research and development that deals with the manipulation of matter at the atomic or subatomic level. It is used in fields such as medicine, energy, aerospace engineering, and more. One of the applications of nanotechnology is in drug delivery. This is a process by which nanoparticles are used to carry and deliver drugs to a specific area in the body. There are several advantages of using nanotechnology for drug delivery, including precise targeting of specific cells, increased drug potency, and lowered toxicity to the cells that are targeted. Nanoparticles can also carry vaccines to cells that might be hard to reach with traditional delivery methods. However, there are some concerns with the use of nanoparticles for drug delivery. Some studies have shown that nanoparticles may contribute to the development of tumors in other parts of the body. There is also growing concern that nanoparticles may have harmful effects on the environment. Despite these potential drawbacks, the use of nanotechnology in drug delivery is still a promising area for future research.
Targeted delivery
Targeted drug delivery is the delivery of a drug to its target site without having an effect on other tissues. Interest in targeted drug delivery has grown drastically due to its potential implications in the treatment of cancers and other chronic diseases. In order to achieve efficient targeted delivery, the designed system must avoid the host's defense mechanisms and circulate to its intended site of action. A number of drug carriers have been studied to effectively target specific tissues, including liposomes, nanogels, and other nanotechnologies.
Controlled-release formulations
Controlled or modified-release formulations alter the rate and timing at which a drug is liberated, in order to produce adequate or sustained drug concentrations. The first controlled-release (CR) formulation that was developed was Dexedrine in the 1950s. This period of time saw more drugs being formulated as CR, as well as the introduction of transdermal patches to allow drugs to slowly absorb through the skin. Since then, countless other CR products have been developed to account for the physiochemical properties of different drugs, such as depot injections for antipsychotics and sex hormones that require dosing once every few months.
Since the late 1990s, most of the research around CR formulations has been focused on implementing nanoparticles to decrease the rate of drug clearance.
Modulated drug release and zero-order drug release
Many scientists worked to create oral formulations that could maintain a constant drug level because of the ability of drug release at a zero-order rate.blood's concentration. However, a few physiological restrictions made it challenging to create such oral formulations. First, because the lower parts of the intestine have a decreased capacity for absorption, the medication absorption typically declines as an oral formulation moves from the stomach to the intestine. The decreased drug amount released from the formulation over time frequently made this condition worse. Phenylpropanolamine HCl release from was the only instance of sustaining consistent blood concentration for roughly 16 hours.
Delivery of biologic drugs
Pharmaceutical preparations containing peptides, proteins, antibodies, genes, or other biologic components often face absorption issues due to their large sizes or electrostatic charges, and may be susceptible to enzymatic degradation once they have entered the body. For these reasons, recent efforts in drug delivery have been focused on methods to avoid these issues through the use of liposomes, nanoparticles, fusion proteins, protein-cage nanoparticles, exploiting routes for the delivery of biologicals that toxins use and many others. Intracellular delivery of macromolecules by chemical carriers is most advanced for RNA, as known from RNA-based COVID-19 vaccines, while proteins have also been delivered into cells in vivo and DNA is routinely delivered in vitro. Among the various routes of administration the oral route is most favored by patients. For most biologic drugs, however, oral bioavailability is too low to reach a therapeutic level. Advanced delivery systems such as formulations containing permeation enhancers or enzyme inhibitors, lipid-based nanocarriers and microneedles will likely enhance oral bioavailability of these drugs sufficiently.
Nanoparticle drug delivery
Drug delivery systems have been around for many years, but there are a few recent applications of drug delivery that warrant
1. Drug delivery to the brain: Many drugs can be harmful when administered systemically; the brain is very sensitive to medications and can easily cause damage if a drug is administered directly into the bloodstream. As new drug formulations are being developed for brain diseases, including Alzheimer's disease and Parkinson's disease, researchers are working on ways to deliver drugs into the brain that do not cause damage to healthy tissue. For example, scientists have developed nanoparticles that can cross the protective blood-brain barrier and deliver drugs directly to the brain.
| Biology and health sciences | General concepts_2 | Health |
487841 | https://en.wikipedia.org/wiki/Cerebrum | Cerebrum | The cerebrum (: cerebra), telencephalon or endbrain is the largest part of the brain, containing the cerebral cortex (of the two cerebral hemispheres) as well as several subcortical structures, including the hippocampus, basal ganglia, and olfactory bulb. In the human brain, the cerebrum is the uppermost region of the central nervous system. The cerebrum develops prenatally from the forebrain (prosencephalon). In mammals, the dorsal telencephalon, or pallium, develops into the cerebral cortex, and the ventral telencephalon, or subpallium, becomes the basal ganglia. The cerebrum is also divided into approximately symmetric left and right cerebral hemispheres.
With the assistance of the cerebellum, the cerebrum controls all voluntary actions in the human body.
Structure
The cerebrum is the largest part of the brain. Depending upon the position of the animal, it lies either in front or on top of the brainstem. In humans, the cerebrum is the largest and best-developed of the five major divisions of the brain.
The cerebrum is made up of the two cerebral hemispheres and their cerebral cortices (the outer layers of grey matter), and the underlying regions of white matter. Its subcortical structures include the hippocampus, basal ganglia and olfactory bulb. The cerebrum consists of two C-shaped cerebral hemispheres, separated from each other by a deep fissure called the longitudinal fissure.
Cerebral cortex
The cerebral cortex, the outer layer of grey matter of the cerebrum, is found only in mammals. In larger mammals, including humans, the surface of the cerebral cortex folds to create gyri (ridges) and sulci (furrows) which increase the surface area.
The cerebral cortex is generally classified into four lobes: the frontal, parietal, occipital and temporal lobes. The lobes are classified based on their overlying neurocranial bones. A smaller lobe is the insular lobe, a part of the cerebral cortex folded deep within the lateral sulcus that separates the temporal lobe from the parietal and frontal lobes, is located within each hemisphere of the mammalian brain.
Cerebral hemispheres
The cerebrum is divided by the medial longitudinal fissure into two cerebral hemispheres, the right and the left. The cerebrum is contralaterally organized, i.e., the right hemisphere controls and processes signals (predominantly) from the left side of the body, while the left hemisphere controls and processes signals (predominantly) from the right side of the body. According to current knowledge, this is due to an axial twist that occurs in the early embryo. There is a strong but not complete bilateral symmetry between the hemispheres, while lateralization tends to increase with increasing brain size. The lateralization of brain function looks at the known and possible differences between the two.
Development
In the developing vertebrate embryo, the neural tube is subdivided into four unseparated sections which then develop further into distinct regions of the central nervous system; these are the prosencephalon (forebrain), the mesencephalon (midbrain) the rhombencephalon (hindbrain) and the spinal cord. The prosencephalon develops further into the telencephalon and the diencephalon. The dorsal telencephalon gives rise to the pallium (cerebral cortex in mammals and reptiles) and the ventral telencephalon generates the basal ganglia. The diencephalon develops into the thalamus and hypothalamus, including the optic vesicles (future retina). The dorsal telencephalon then forms two lateral telencephalic vesicles, separated by the midline, which develop into the left and right cerebral hemispheres. Birds and fish have a dorsal telencephalon, like all vertebrates, but it is generally unlayered and therefore not considered a cerebral cortex. Only a layered cytoarchitecture can be considered a cortex.
Functions
Note: As the cerebrum is a gross division with many subdivisions and sub-regions, it is important to state that this section lists only functions that the cerebrum as a whole serves. (See main articles on cerebral cortex and basal ganglia for more information.) The cerebrum is a major part of the brain, controlling emotions, hearing, vision, personality and much more. It controls all precision of voluntary actions, and it functions as the center of sensory perception, memory, thoughts and judgement; the cerebrum also functions as the center of voluntary motor activities.
Motor functions
Upper motor neurons in the primary motor cortex send their axons to the brainstem and spinal cord to synapse on the lower motor neurons, which innervate the muscles. Damage to motor areas of the cortex can lead to certain types of motor neuron disease. This kind of damage results in loss of muscular power and precision rather than total paralysis.
Sensory processing
The primary sensory areas of the cerebral cortex receive and process visual, auditory, somatosensory, gustatory, and olfactory information. Together with association cortical areas, these brain regions synthesize sensory information into our perceptions of the world.
Olfaction
The olfactory bulb, responsible for the sense of smell, takes up a large area of the cerebrum in most vertebrates. However, in humans, this part of the brain is much smaller and lies underneath the frontal lobe. The olfactory sensory system is unique since the neurons in the olfactory bulb send their axons directly to the olfactory cortex, rather than to the thalamus first. Olfaction is also the only sense that is represented by the ipsilateral side of the brain. Damage to the olfactory bulb results in a loss of olfaction (the sense of smell).
Language and communication
Speech and language are mainly attributed to parts of the cerebral cortex. Motor portions of language are attributed to Broca's area within the frontal lobe. Speech comprehension is attributed to Wernicke's area, at the temporal-parietal lobe junction. These two regions are interconnected by a large white matter tract, the arcuate fasciculus. Damage to the Broca's area results in expressive aphasia (non-fluent aphasia) while damage to Wernicke's area results in receptive aphasia (also called fluent aphasia).
Learning and memory
Explicit or declarative (factual) memory formation is attributed to the hippocampus and associated regions of the medial temporal lobe. This association was originally described after a patient known as HM had both his left and right hippocampus surgically removed to treat chronic [temporal lobe epilepsy]. After surgery, HM had anterograde amnesia, or the inability to form new memories.
Implicit or procedural memory, such as complex motor behaviors, involves the basal ganglia.
Short-term or working memory involves association areas of the cortex, especially the dorsolateral prefrontal cortex, as well as the hippocampus.
Other animals
In the most primitive vertebrates, the hagfishes and lampreys, the cerebrum is a relatively simple structure receiving nerve impulses from the olfactory bulb. In cartilaginous and lobe-finned fishes and also in amphibians, a more complex structure is present, with the cerebrum being divided into three distinct regions. The lowermost (or ventral) region forms the basal nuclei, and contains fibres connecting the rest of the cerebrum to the thalamus. Above this, and forming the lateral part of the cerebrum, is the paleopallium, while the uppermost (or dorsal) part is referred to as the archipallium. The cerebrum remains largely devoted to olfactory sensation in these animals, in contrast to its much wider range of functions in amniotes.
In ray-finned fishes, the structure is somewhat different. The inner surfaces of the lateral and ventral regions of the cerebrum bulge up into the ventricles; these include both the basal nuclei and the various parts of the pallium and may be complex in structure, especially in teleosts. The dorsal surface of the cerebrum is membranous, and does not contain any nervous tissue.
In the amniotes, the cerebrum becomes increasingly large and complex. In reptiles, the paleopallium is much larger than in amphibians and its growth has pushed the basal nuclei into the central regions of the cerebrum. As in the lower vertebrates, the grey matter is generally located beneath the white matter, but in some reptiles, it spreads out to the surface to form a primitive cortex, especially in the anterior part of the brain.
In mammals, this development proceeds further, so that the cortex covers almost the whole of the cerebral hemispheres, especially in more developed species, such as the primates. The paleopallium is pushed to the ventral surface of the brain, where it becomes the olfactory lobes, while the archipallium becomes rolled over at the medial dorsal edge to form the hippocampus. In placental mammals, a corpus callosum also develops, further connecting the two hemispheres. The complex convolutions of the cerebral surface (see gyrus, gyrification) are also found only in higher mammals. Although some large mammals (such as elephants) have particularly large cerebra, dolphins are the only species (other than humans) to have cerebra accounting for as much as 2 percent of their body weight.
The cerebra of birds are similarly enlarged to those of mammals, by comparison with reptiles. The increased size of bird brains was classically attributed to enlarged basal ganglia, with the other areas remaining primitive, but this view has been largely abandoned. Birds appear to have undergone an alternate process of encephalization, as they diverged from the other archosaurs, with few clear parallels to that experienced by mammals and their therapsid ancestors.
Additional images
| Biology and health sciences | Nervous system | null |
487859 | https://en.wikipedia.org/wiki/Directional%20derivative | Directional derivative | In multivariable calculus, the directional derivative measures the rate at which a function changes in a particular direction at a given point.
The directional derivative of a multivariable differentiable (scalar) function along a given vector v at a given point x intuitively represents the instantaneous rate of change of the function, moving through x with a direction specified by v.
The directional derivative of a scalar function f with respect to a vector v at a point (e.g., position) x may be denoted by any of the following:
It therefore generalizes the notion of a partial derivative, in which the rate of change is taken along one of the curvilinear coordinate curves, all other coordinates being constant.
The directional derivative is a special case of the Gateaux derivative.
Definition
The directional derivative of a scalar function
along a vector
is the function defined by the limit
This definition is valid in a broad range of contexts, for example where the norm of a vector (and hence a unit vector) is undefined.
For differentiable functions
If the function f is differentiable at x, then the directional derivative exists along any unit vector v at x, and one has
where the on the right denotes the gradient, is the dot product and v is a unit vector. This follows from defining a path and using the definition of the derivative as a limit which can be calculated along this path to get:
Intuitively, the directional derivative of f at a point x represents the rate of change of f, in the direction of v.
Using only direction of vector
In a Euclidean space, some authors define the directional derivative to be with respect to an arbitrary nonzero vector v after normalization, thus being independent of its magnitude and depending only on its direction.
This definition gives the rate of increase of per unit of distance moved in the direction given by . In this case, one has
or in case f is differentiable at x,
Restriction to a unit vector
In the context of a function on a Euclidean space, some texts restrict the vector v to being a unit vector. With this restriction, both the above definitions are equivalent.
Properties
Many of the familiar properties of the ordinary derivative hold for the directional derivative. These include, for any functions f and g defined in a neighborhood of, and differentiable at, p:
sum rule:
constant factor rule: For any constant c,
product rule (or Leibniz's rule):
chain rule: If g is differentiable at p and h is differentiable at g(p), then
In differential geometry
Let be a differentiable manifold and a point of . Suppose that is a function defined in a neighborhood of , and differentiable at . If is a tangent vector to at , then the directional derivative of along , denoted variously as (see Exterior derivative), (see Covariant derivative), (see Lie derivative), or (see ), can be defined as follows. Let be a differentiable curve with and . Then the directional derivative is defined by
This definition can be proven independent of the choice of , provided is selected in the prescribed manner so that and .
The Lie derivative
The Lie derivative of a vector field along a vector field is given by the difference of two directional derivatives (with vanishing torsion):
In particular, for a scalar field , the Lie derivative reduces to the standard directional derivative:
The Riemann tensor
Directional derivatives are often used in introductory derivations of the Riemann curvature tensor. Consider a curved rectangle with an infinitesimal vector along one edge and along the other. We translate a covector along then and then subtract the translation along and then . Instead of building the directional derivative using partial derivatives, we use the covariant derivative. The translation operator for is thus
and for ,
The difference between the two paths is then
It can be argued that the noncommutativity of the covariant derivatives measures the curvature of the manifold:
where is the Riemann curvature tensor and the sign depends on the sign convention of the author.
In group theory
Translations
In the Poincaré algebra, we can define an infinitesimal translation operator P as
(the i ensures that P is a self-adjoint operator) For a finite displacement λ, the unitary Hilbert space representation for translations is
By using the above definition of the infinitesimal translation operator, we see that the finite translation operator is an exponentiated directional derivative:
This is a translation operator in the sense that it acts on multivariable functions f(x) as
Rotations
The rotation operator also contains a directional derivative. The rotation operator for an angle θ, i.e. by an amount θ = |θ| about an axis parallel to is
Here L is the vector operator that generates SO(3):
It may be shown geometrically that an infinitesimal right-handed rotation changes the position vector x by
So we would expect under infinitesimal rotation:
It follows that
Following the same exponentiation procedure as above, we arrive at the rotation operator in the position basis, which is an exponentiated directional derivative:
Normal derivative
A normal derivative is a directional derivative taken in the direction normal (that is, orthogonal) to some surface in space, or more generally along a normal vector field orthogonal to some hypersurface. See for example Neumann boundary condition. If the normal direction is denoted by , then the normal derivative of a function f is sometimes denoted as . In other notations,
In the continuum mechanics of solids
Several important results in continuum mechanics require the derivatives of vectors with respect to vectors and of tensors with respect to vectors and tensors. The directional directive provides a systematic way of finding these derivatives.
| Mathematics | Multivariable and vector calculus | null |
487896 | https://en.wikipedia.org/wiki/Stickleback | Stickleback | The sticklebacks are a family of ray-finned fishes, the Gasterosteidae which have a Holarctic distribution in fresh, brackish and marine waters. They were thought to be related to the pipefish and seahorses but are now thought to be more closely related to the eelpouts and sculpins.
Taxonomy
The stickleback family, Gasterosteidae, was first proposed as a family by the French zoologist Charles Lucien Bonaparte in 1831. It was long thought that the sticklebacks and their relatives made up a suborder, the Gasterosteoidei, of the order Gasterostiformes with the sea horses and pipefishes making up the suborder Syngnathoidei. More recent phylogenetic work has shown that the Gaterosteoidei is more closely related to the Zoarcoidei and the Cottoidei, which means that this taxon would belong in the order Scorpaeniformes. but in other phylogenetic classifications it is treated as the infraorder Gasterosteales within the suborder Cottoidei or as a sister clade to the Zoarcales in the order Zoarciformes.
FishBase recognises 16 species in the family, grouped in five genera. However, several of the species have a number of recognised subspecies, and the taxonomy of the family is thought to be in need of revision.
Genera
The family Gasterosteidae includes the following genera:
Apeltes DeKay, 1842
Culaea Whitley, 1950
Gasterosteus Linnaeus, 1758
Pungitius d'Annone, 1760
Spinachia Cuvier, 1816
Description
Sticklebacks are endemic to the temperate zone and are most commonly found in the ocean, but some can be found in fresh water. The freshwater taxa were trapped in Europe, Asia, and North America after the Ice Age 10,000–20,000 years ago, and have evolved features different from those of the marine species.
Sticklebacks are carnivorous, feeding on small animals such as insects, crustaceans and fish larvae.
Sticklebacks are characterised by the presence of strong and clearly isolated spines in their dorsal fins. An unusual feature of sticklebacks is that they have no scales, although some species have bony armour plates.
Sizes
The maximum size of the best-known species, the three-spined stickleback (Gasterosteus aculeatus), is about 4 inches, but few of them are more than 3 inches long. They mature sexually at a length of about 2 inches. Most other stickleback species are roughly similar in size or somewhat smaller. The only exception is the far larger fifteen-spined stickleback (Spinachia spinachia), which can reach 22 cm (approx. 8.8 inches). Body form varies with habitat: sticklebacks in shallow lakes have developed a deep body specialized to enable feeding on benthic invertebrates, whilst those in deep oligotrophic lakes have adapted to feed on plankton and have a more slimlined body.
Personality
Research has shown that Sticklebacks display distinct personality traits, specifically in the area of taking a risk, and, can be considered bold or shy. These personality traits were determined to directly influence if they would lead, and if discouraged, attempt to lead again.
Mating
All stickleback species show similar, unusual mating behaviour. Freshwater males develop a red colouration, and although this may be seen in oceanic and benthic species these tend to remain dull-coloured. The male then constructs a nest from weeds held together by spiggin, a kidney secretion, then attract females to the nest. Females signal their readiness to mate with solitary rather than shoaling behaviour, a head-up posture; their bellies are also obviously distended with eggs. Courtship typically involves a zig-zag 'dance' where the male approaches the female in an erratic side-to-side pattern, and dorsal pricking of the female's abdomen. A female lays her eggs inside the nest, where the male fertilises them. The male then guards the eggs until they hatch 7–14 days later (depending on temperature), and may continue to guard the fry after they hatch. This large investment in both the nesting site and guarding of the eggs limits the number of females a male can mate with however males spawn multiple times. This introduces the ability for selection to favor male mate choice. Some males die following spawning.
Mating choice
Typically, the sex with the greatest parental investment has the strongest mate preferences. Stickleback species exhibit mutual mate choice in which both the male and female have strong mate preferences. This is due in part to the strong parental investment on behalf of the male in guarding the eggs.
Female mate choice
Female sticklebacks show a strong preference to male stickleback with bright red coloration under their throats. Females mate both more often with males with brighter red coloration and give on average, larger eggs to be fertilized by these males. This preference has led to brighter red coloring. This association is possible because the red coloration can only be produced by males that are free of parasites. This is referred to in the Hamilton-Zuk hypothesis.
However, there is also evidence that attractive male red coloration may be a faulty signal of male quality. Male sticklebacks that are more attractive to females due to carotenoid colorants may under-allocate carotenoids to their germline cells. Since carotinoids are beneficial antioxidants, their under-allocation to germline cells can lead to increased oxidative DNA damage to these cells. Therefore, female sticklebacks may risk fertility and the viability of their offspring by choosing redder, but more deteriorated partners with reduced sperm quality.
Female mate choice has also been seen to be condition dependent. Females are almost always the more choosy sex in most species. Female sticklebacks though, have been found to be less choosy of mates when in poor physical condition and inversely, more choosy in good condition.
Male mate choice
In some species, such as the three-spined stickleback, the large investment in both nesting site and guarding of eggs by males limits the number of females a male can mate with. This introduces the ability for selection to favor male mate choice. Male mate choice is rarely studied or observed in many species but multiple studies have confirmed male mate choice within stickleback species. Males show a choosiness similar to females as to what female they are willing to court and mate. Male sticklebacks have been observed to show preference towards female sticklebacks that are larger and longer. This is believed to be because larger females on average produce larger eggs, which leads to a greater offspring survival and fitness. In addition, male sticklebacks have also been observed to prefer females with more distended or bloated stomachs. The benefits of this is also due to larger eggs and thus offspring survival and fitness
Inbreeding avoidance
Female three-spined sticklebacks adjust their courting behaviour to the risk of inbreeding. When gravid females are given the choice between a courting unfamiliar non-sibling and a familiar brother, they prefer to mate with the non-sibling and thus avoid the disadvantages that accompany incest. Eggs from inbred matings compared to eggs from outbred matings have a lower rate of fertilization and hatching, and fewer progeny survive to reproductive age.
Use in science
Niko Tinbergen's studies of the behaviour of this fish were important in the early development of ethology as an example of a fixed action pattern. More recently, the fish have become a favourite system for studying the molecular genetics of evolutionary change in wild populations and a powerful "supermodel" for combining evolutionary studies at molecular, developmental, population genetic, and ecological levels. The nearly complete genome sequence of a reference freshwater stickleback was described in 2012, along with set of genetic variants commonly found in 21 marine and freshwater populations around the world. Some variants, and several chromosome inversions, consistently distinguish marine and freshwater populations, helping identify a genome-wide set of changes contributing to repeated adaptation of sticklebacks to marine and freshwater environments. The adaptations seen in oceanic threespine sticklebacks make them an ideal organism for the study of parallel evolution.
In culture
There is a sculpture in Kronstadt dedicated to stickleback, which saved thousands of city residents from starvation during the Leningrad Siege of World War II.
| Biology and health sciences | Fishes | null |
488341 | https://en.wikipedia.org/wiki/Homo | Homo | Homo () is a genus of great ape (family Hominidae) that emerged from the genus Australopithecus and encompasses only a single extant species, Homo sapiens (modern humans), along with a number of extinct species (collectively called archaic humans) classified as either ancestral or closely related to modern humans; these include Homo erectus and Homo neanderthalensis. The oldest member of the genus is Homo habilis, with records of just over 2 million years ago. Homo, together with the genus Paranthropus, is probably most closely related to the species Australopithecus africanus within Australopithecus. The closest living relatives of Homo are of the genus Pan (chimpanzees and bonobos), with the ancestors of Pan and Homo estimated to have diverged around 5.7-11 million years ago during the Late Miocene.
H. erectus appeared about 2 million years ago and spread throughout Africa (debatably as another species called Homo ergaster) and Eurasia in several migrations. The species was adaptive and successful, and persisted for more than a million years before gradually diverging into new species around 500,000 years ago.
Anatomically modern humans (H. sapiens) emerged close to 300,000 to 200,000 years ago in Africa, and H. neanderthalensis emerged around the same time in Europe and Western Asia. H. sapiens dispersed from Africa in several waves, from possibly as early as 250,000 years ago, and certainly by 130,000 years ago, with the so-called Southern Dispersal, beginning about 70–50,000 years ago, leading to the lasting colonisation of Eurasia and Oceania by 50,000 years ago. H. sapiens met and interbred with archaic humans in Africa and in Eurasia. Separate archaic (non-sapiens) human species including Neanderthals are thought to have survived until around 40,000 years ago.
Names and taxonomy
The Latin noun homō (genitive hominis) means "human being" or "man" in the generic sense of "human being, mankind". The binomial name Homo sapiens was coined by Carl Linnaeus (1758). Names for other species of the genus were introduced from the second half of the 19th century (H. neanderthalensis 1864, H. erectus 1892).
The genus Homo has not been strictly defined, even today. Since the early human fossil record began to slowly emerge from the earth, the boundaries and definitions of the genus have been poorly defined and constantly in flux. Because there was no reason to think it would ever have any additional members, Carl Linnaeus did not even bother to define Homo when he first created it for humans in the 18th century. The discovery of Neanderthal brought the first addition.
The genus Homo was given its taxonomic name to suggest that its member species can be classified as human. And, over the decades of the 20th century, fossil finds of pre-human and early human species from late Miocene and early Pliocene times produced a rich mix for debating classifications. There is continuing debate on delineating Homo from Australopithecus—or, indeed, delineating Homo from Pan. Even so, classifying the fossils of Homo coincides with evidence of: (1) competent human bipedalism in Homo habilis inherited from the earlier Australopithecus of more than four million years ago, as demonstrated by the Laetoli footprints; and (2) human tool culture having begun by 2.5 million years ago to 3 million years ago.
From the late-19th to mid-20th centuries, a number of new taxonomic names, including new generic names, were proposed for early human fossils; most have since been merged with Homo in recognition that Homo erectus was a single species with a large geographic spread of early migrations. Many such names are now regarded as "synonyms" with Homo, including Pithecanthropus, Protanthropus, Sinanthropus, Cyphanthropus, Africanthropus, Telanthropus, Atlanthropus, and Tchadanthropus.
Classifying the genus Homo into species and subspecies is subject to incomplete information and remains poorly done. This has led to using common names ("Neanderthal" and "Denisovan"), even in scientific papers, to avoid trinomial names or the ambiguity of classifying groups as incertae sedis (uncertain placement)—for example, H. neanderthalensis vs. H. sapiens neanderthalensis, or H. georgicus vs. H. erectus georgicus. Some recently extinct species in the genus have been discovered only lately and do not as yet have consensus binomial names (see Denisova hominin). Since the beginning of the Holocene, it is likely that Homo sapiens (anatomically modern humans) has been the only extant species of Homo.
John Edward Gray (1825) was an early advocate of classifying taxa by designating tribes and families. Wood and Richmond (2000) proposed that Hominini ("hominins") be designated as a tribe that comprised all species of early humans and pre-humans ancestral to humans back to after the chimpanzee–human last common ancestor, and that Hominina be designated a subtribe of Hominini to include only the genus Homo — that is, not including the earlier upright walking hominins of the Pliocene such as Australopithecus, Orrorin tugenensis, Ardipithecus, or Sahelanthropus. Designations alternative to Hominina existed, or were offered: Australopithecinae (Gregory & Hellman 1939) and Preanthropinae (Cela-Conde & Altaba 2002); and later, Cela-Conde and Ayala (2003) proposed that the four genera Australopithecus, Ardipithecus, Praeanthropus, and Sahelanthropus be grouped with Homo within Hominini (sans Pan).
Evolution
Australopithecus and the appearance of Homo
Several species, including Australopithecus garhi, Australopithecus sediba, Australopithecus africanus, and Australopithecus afarensis, have been proposed as the ancestor or sister of the Homo lineage. These species have morphological features that align them with Homo, but there is no consensus as to which gave rise to Homo.
Especially since the 2010s, the delineation of Homo in Australopithecus has become more contentious. Traditionally, the advent of Homo has been taken to coincide with the first use of stone tools (the Oldowan industry), and thus by definition with the beginning of the Lower Palaeolithic. But in 2010, evidence was presented that seems to attribute the use of stone tools to Australopithecus afarensis around 3.3 million years ago, close to a million years before the first appearance of Homo. LD 350-1, a fossil mandible fragment dated to 2.8 Mya, discovered in 2013 in Afar, Ethiopia, was described as combining "primitive traits seen in early Australopithecus with derived morphology observed in later Homo. Some authors would push the development of Homo close to or even past 3 Mya. This finds support in a recent phylogenetic study in hominins that by using morphological, molecular and radiometric information, dates the emergence of Homo at 3.3 Ma (4.30 – 2.56 Ma). Others have voiced doubt as to whether Homo habilis should be included in Homo, proposing an origin of Homo with Homo erectus at roughly 1.9 Mya instead.
The most salient physiological development between the earlier australopithecine species and Homo is the increase in endocranial volume (ECV), from about in A. garhi to in H. habilis and further to in H. erectus, in H. heidelbergensis and up to in H. neanderthalensis. However, a steady rise in cranial capacity is observed already in Autralopithecina and does not terminate after the emergence of Homo, so that it does not serve as an objective criterion to define the emergence of the genus.
Homo habilis
Homo habilis emerged about 2.1 Mya. Already before 2010, there were suggestions that H. habilis should not be placed in the genus Homo but rather in Australopithecus. The main reason to include H. habilis in Homo, its undisputed tool use, has become obsolete with the discovery of Australopithecus tool use at least a million years before H. habilis. Furthermore, H. habilis was long thought to be the ancestor of the more gracile Homo ergaster (Homo erectus). In 2007, it was discovered that H. habilis and H. erectus coexisted for a considerable time, suggesting that H. erectus is not immediately derived from H. habilis but instead from a common ancestor. With the publication of Dmanisi skull 5 in 2013, it has become less certain that Asian H. erectus is a descendant of African H. ergaster which was in turn derived from H. habilis. Instead, H. ergaster and H. erectus appear to be variants of the same species, which may have originated in either Africa or Asia and widely dispersed throughout Eurasia (including Europe, Indonesia, China) by 0.5 Mya.
Homo erectus
Homo erectus has often been assumed to have developed anagenetically from H. habilis from about 2 million years ago. This scenario was strengthened with the discovery of Homo erectus georgicus, early specimens of H. erectus found in the Caucasus, which seemed to exhibit with H. habilis. As the earliest evidence for H. erectus was found outside of Africa, it was considered plausible that H. erectus developed in Eurasia and then migrated back to Africa. Based on fossils from the Koobi Fora Formation, east of Lake Turkana in Kenya, Spoor et al. (2007) argued that H. habilis may have survived beyond the emergence of H. erectus, so that the evolution of H. erectus would not have been anagenetically, and H. erectus would have existed alongside H. habilis for about half a million years (), during the early Calabrian. On 31 August 2023, researchers reported, based on genetic studies, that a human ancestor population bottleneck (from a possible 100,000 to 1000 individuals) occurred "around 930,000 and 813,000 years ago ... lasted for about 117,000 years and brought human ancestors close to extinction."
Weiss (1984) estimated that there have been about 44 billion (short scale) members of the genus Homo from its origins to the evolution of H. erectus, about 56 billion individuals from H. erectus to the Neolithic, and another 51 billion individuals since the Neolithic. This provides the opportunity for an immense amount of new mutational variation to have arisen during human evolution.
A separate South African species Homo gautengensis has been postulated as contemporary with H. erectus in 2010.
Phylogeny
A taxonomy of Homo within the great apes is assessed as follows, with Paranthropus and Homo emerging within Australopithecus (shown here cladistically granting Paranthropus, Kenyanthropus, and Homo). The exact phylogeny within Australopithecus is still highly controversial. Approximate radiation dates of daughter clades are shown in millions of years ago (Mya). Sahelanthropus and Orrorin, possibly sisters to Australopithecus, are not shown here. The naming of groupings is sometimes muddled as often certain groupings are presumed before any cladistic analysis is performed.
Cladogram based on Dembo et al. (2016):
Cladogram based on Feng et al. (2024):
Several of the Homo lineages appear to have surviving progeny through introgression into other lines. Genetic evidence indicates an archaic lineage separating from the other human lineages 1.5 million years ago, perhaps H. erectus, may have interbred into the Denisovans about 55,000 years ago. Fossil evidence shows H. erectus s.s. survived at least until 117,000 yrs ago, and the even more basal H. floresiensis survived until 50,000 years ago. A 1.5-million-year H. erectus-like lineage appears to have made its way into modern humans through the Denisovans and specifically into the Papuans and aboriginal Australians. The genomes of non-sub-Saharan African humans show what appear to be numerous independent introgression events involving Neanderthal and in some cases also Denisovans around 45,000 years ago. The genetic structure of some sub-Saharan African groups seems to be indicative of introgression from a west Eurasian population some 3,000 years ago.
Some evidence suggests that Australopithecus sediba could be moved to the genus Homo, or placed in its own genus, due to its position with respect to e.g. H. habilis and H. floresiensis.
Dispersal
By about 1.8 million years ago, H. erectus is present in both East Africa (H. ergaster) and in Western Asia (H. georgicus). The ancestors of Indonesian H. floresiensis may have left Africa even earlier.
Homo erectus and related or derived archaic human species over the next 1.5 million years spread throughout Africa and Eurasia (see: Recent African origin of modern humans). Europe is reached by about 0.5 Mya by Homo heidelbergensis.
Homo neanderthalensis and H. sapiens develop after about 300 kya. Homo naledi is present in Southern Africa by 300 kya.
H. sapiens soon after its first emergence spread throughout Africa, and to Western Asia in several waves, possibly as early as 250 kya, and certainly by 130 kya. In July 2019, anthropologists reported the discovery of 210,000 year old remains of a H. sapiens and 170,000 year old remains of a H. neanderthalensis in Apidima Cave, Peloponnese, Greece, more than 150,000 years older than previous H. sapiens finds in Europe.
Most notable is the Southern Dispersal of H. sapiens around 60 kya, which led to the lasting peopling of Oceania and Eurasia by anatomically modern humans. H. sapiens interbred with archaic humans both in Africa and in Eurasia, in Eurasia notably with Neanderthals and Denisovans.
Among extant populations of H. sapiens, the deepest temporal division is found in the San people of Southern Africa, estimated at close to 130,000 years, or possibly more than 300,000 years ago. Temporal division among non-Africans is of the order of 60,000 years in the case of Australo-Melanesians. Division of Europeans and East Asians is of the order of 50,000 years, with repeated and significant admixture events throughout Eurasia during the Holocene.
Archaic human species may have survived until the beginning of the Holocene, although they were mostly extinct or absorbed by the expanding H. sapiens populations by 40 kya (Neanderthal extinction).
List of lineages
The species status of H. rudolfensis, H. ergaster, H. georgicus, H. antecessor, H. cepranensis, H. rhodesiensis, H. neanderthalensis, Denisova hominin, and H. floresiensis remain under debate. H. heidelbergensis and H. neanderthalensis are closely related to each other and have been considered to be subspecies of H. sapiens.
There has historically been a trend to postulate new human species based on as little as an individual fossil. A "minimalist" approach to human taxonomy recognizes at most three species, H. habilis (2.1–1.5 Mya, membership in Homo questionable), H. erectus (1.8–0.1 Mya, including the majority of the age of the genus, and the majority of archaic varieties as subspecies, including H. heidelbergensis as a late or transitional variety) and Homo sapiens (300 kya to present, including H. neanderthalensis and other varieties as subspecies). Consistent definitions and methodology of species delineation are not generally agreed upon in anthropology or paleontology. Indeed, speciating populations of mammals can typically interbreed for several million years after they begin to genetically diverge, so all contemporary "species" in the genus Homo would potentially have been able to interbreed at the time, and introgression from beyond the genus Homo can not a priori be ruled out. It has been suggested that H. naledi may have been a hybrid with a late surviving Australipith (taken to mean beyond Homo, ed.), despite the fact that these lineages generally are regarded as long extinct. As discussed above, many introgressions have occurred between lineages, with evidence of introgression after separation of 1.5 million years.
| Biology and health sciences | Homo | Biology |
488361 | https://en.wikipedia.org/wiki/Mining%20engineering | Mining engineering | Mining in the engineering discipline is the extraction of minerals from the ground. Mining engineering is associated with many other disciplines, such as mineral processing, exploration, excavation, geology, metallurgy, geotechnical engineering and surveying. A mining engineer may manage any phase of mining operations, from exploration and discovery of the mineral resources, through feasibility study, mine design, development of plans, production and operations to mine closure.
History of mining engineering
From prehistoric times to the present, mining has played a significant role in the existence of the human race. Since the beginning of civilization, people have used stone and ceramics and, later, metals found on or close to the Earth's surface. These were used to manufacture early tools and weapons. For example, high-quality flint found in northern France and southern England were used to set fire and break rock. Flint mines have been found in chalk areas where seams of the stone were followed underground by shafts and galleries. The oldest known mine on the archaeological record is the "Lion Cave" in Eswatini. At this site, which radiocarbon dating indicates to be about 43,000 years old, paleolithic humans mined mineral hematite, which contained iron and was ground to produce the red pigment ochre.
The ancient Romans were innovators of mining engineering. They developed large-scale mining methods, such as the use of large volumes of water brought to the minehead by aqueducts for hydraulic mining. The exposed rock was then attacked by fire-setting, where fires were used to heat the rock, which would be quenched with a stream of water. The thermal shock cracked the rock, enabling it to be removed. In some mines, the Romans utilized water-powered machinery such as reverse overshot water-wheels. These were used extensively in the copper mines at Rio Tinto in Spain, where one sequence comprised 16 such wheels arranged in pairs, lifting water about .
Black powder was first used in mining in Banská Štiavnica, Kingdom of Hungary (present-day Slovakia) in 1627. This allowed blasting of rock and earth to loosen and reveal ore veins, which was much faster than fire-setting. The Industrial Revolution saw further advances in mining technologies, including improved explosives and steam-powered pumps, lifts, and drills.
Education
Becoming an accredited mining engineer requires a university or college degree. Training includes a Bachelor of Engineering (B.Eng. or B.E.), Bachelor of Science (B.Sc. or B.S.), Bachelor of Technology (B.Tech.) or Bachelor of Applied Science (B.A.Sc.) in mining engineering. Depending on the country and jurisdiction, to be licensed as a mining engineer may require a Master of Engineering (M.Eng.), Master of Science (M.Sc or M.S.) or Master of Applied Science (M.A.Sc.) degree.
Some mining engineers who have come from other disciplines, primarily from engineering fields (e.g.: mechanical, civil, electrical, geomatics or environmental engineering) or from science fields (e.g.: geology, geophysics, physics, geomatics, earth science, or mathematics), typically completing a graduate degree such as M.Eng, M.S., M.Sc. or M.A.Sc. in mining engineering after graduating from a different quantitative undergraduate program.
The fundamental subjects of mining engineering study usually include:
mathematics; calculus, algebra, numerical analysis, statistics
geoscience; geochemistry, geophysics, mineralogy, geomatics
mechanics; rock mechanics, soil Mechanics, geomechanics
thermodynamics; heat transfer, mass transfer
hydrogeology
fluid mechanics; fluid statics, fluid dynamics
Geostatistics; spatial analysis
control engineering; control theory, instrumentation
surface mining; open-pit mining
underground mining (soft rock)
underground mining (hard rock)
computing; DATAMINE, MATLAB, Maptek (Vulcan), Golden Software (Surfer), MicroStation, Carlson
drilling and blasting
solid mechanics; fracture mechanics
In the United States, about 14 universities offer a B.S. degree in mining and mineral engineering. The top rated universities include West Virginia University, South Dakota School of Mines and Technology, Virginia Tech, the University of Kentucky, the University of Arizona, Montana Tech, and Colorado School of Mines. Most of these universities offer M.S. and Ph.D. degrees.
In Canada, there are 19 undergraduate degree programs in mining engineering or equivalent. McGill University Faculty of Engineering offers both undergraduate (B.Sc., B.Eng.) and graduate (M.Sc., Ph.D.) degrees in Mining Engineering. and the University of British Columbia in Vancouver offers a Bachelor of Applied Science (B.A.Sc.) in Mining Engineering and also graduate degrees (M.A.Sc. or M.Eng and Ph.D.) in Mining Engineering.
In Europe, most programs are integrated (B.S. plus M.S. into one) after the Bologna Process and take five years to complete. In Portugal, the University of Porto offers an M.Eng. in Mining and Geo-Environmental Engineering and in Spain the Technical University of Madrid offers degrees in Mining Engineering with tracks in Mining Technology, Mining Operations, Fuels and Explosives, Metallurgy. In the United Kingdom, The Camborne School of Mines offers a wide choice of BEng and MEng degrees in Mining engineering and other Mining related disciplines. This is done through the University of Exeter. In Romania, the University of Petroșani (formerly known as the Petroşani Institute of Mines, or rarely as the Petroşani Institute of Coal) is the only university that offers a degree in Mining Engineering, Mining Surveying or Underground Mining Constructions, albeit, after the closure of Jiu Valley coal mines, those degrees had fallen out of interest for most high-school graduates.
In South Africa, leading institutions include the University of Pretoria, offering a 4-year Bachelor of Engineering (B.Eng in Mining Engineering) as well as post-graduate studies in various specialty fields such as rock engineering and numerical modelling, explosives engineering, ventilation engineering, underground mining methods and mine design; and the University of the Witwatersrand offering a 4-year Bachelor of Science in Engineering (B.Sc.(Eng.)) in Mining Engineering as well as graduate programs (M.Sc.(Eng.) and Ph.D.) in Mining Engineering.
Some mining engineers go on to pursue Doctorate degree programs such as Doctor of Philosophy (Ph.D., DPhil), Doctor of Engineering (D.Eng., Eng.D.). These programs involve a significant original research component and are usually seen as entry points into academia.
In the Russian Federation, 85 universities across all federal districts are training specialists for the mineral resource sector. 36 universities are training specialists for extracting and processing solid minerals (mining). 49 are training specialists for extracting, primary processing, and transporting liquid and gaseous minerals (oil and gas). 37 are training specialists for geological exploration (applied geology, geological exploration). Among the universities that train specialists for the mineral resource sector, 7 are federal universities, and 13 are national research universities of Russia. Personnel training for the mineral resource sector in Russian universities is currently carried out in the following main specializations of training (specialist's degree): "Applied Geology" with the qualification of mining engineer (5 years of training); "Geological Exploration" with the qualification of mining engineer (5 years of training); "Mining" with the qualification of mining engineer (5.5 years of training); "Physical Processes in Mining or Oil and Gas Production" with the qualification of mining engineer (5.5 years of training); "Oil and Gas Engineering and Technologies" with the qualification of mining engineer (5.5 years of training). Universities develop and implement the main professional educational programs of higher education in the directions and specializations of training by forming their profile (name of the program). For example, within the framework of the specialization "Mining", universities often adhere to the classical names of the programs "Open-pit mining", "Underground mining of mineral deposits", "Surveying", "Mineral enrichment", "Mining machines", "Technological safety and mine rescue", "Mine and underground construction", "Blasting work", "Electrification of the mining industry", etc. In the last ten years, under the influence of various factors, new names of programs have begun to appear, such as: "Mining and geological information systems", "Mining ecology", etc. Thus, universities, using their freedom to form new training programs for specialists, can look to the future and try to foresee new professions of mining engineers. After the specialist's degree, you can immediately enrol in postgraduate school (analogue of Doctorate degree programs, four years of training).
Salary and statistics
Mining salaries are usually determined by the level of skill required, where the position is, and what kind of organization the engineer works for.
Mining engineers in India earn relatively high salaries in comparison to many other professions, with an average salary of $15,250 . However, in comparison to mining engineer salaries in other regions, such as Canada, the United States, Australia, and the United Kingdom, Indian salaries are low. In the United States, there are an estimated 6,150 employed mining engineers, with a mean yearly wage of US$103,710.
Pre-mining
As there is considerable capital expenditure required for mining operations, an array of pre-mining activities are normally carried out to assess whether a mining operation would be worthwhile.
Mineral exploration is the process of locating minerals and assessing their concentrations (grade) and quantities (tonnage), to determine if they are commercially viable ores for mining. Mineral exploration is much more intensive, organized, involved, and professional than mineral prospecting – though it frequently utilizes services exploration, enlisting geologists and surveyors in the necessary pre-feasibility study of the possible mining operation. Mineral exploration and estimation of the reserve can determine the profitability conditions and advocate the form and type of mining required.
Mineral discovery
Mineral discovery can be made from research of mineral maps, academic geological reports, or government geological reports. Other sources of information include property assays and local word of mouth. Mineral research usually includes sampling and analysing sediments, soil, and drill cores. Soil sampling and analysis is one of the most popular mineral exploration tools. Other common tools include satellite and aerial surveys or airborne geophysics, including magneto-metric and gamma-spectrometric maps. Unless the mineral exploration is done on public property, the owners of the property may play a significant role in the exploration process and might be the original discoverers of the mineral deposit.
Mineral determination
After a prospective mineral is located, the mining geologist and engineer determine the ore properties. This may involve chemical analysis of the ore to determine the sample's composition. Once the mineral properties are identified, the next step is determining the quantity of the ore. This involves determining the extent of the deposit and the purity of the ore. The geologist drills additional core samples to find the limits of the deposit or seam and estimates the quantity of valuable material present.
Feasibility study
Once the mineral identification and reserve amount are reasonably determined, the next step is to determine the feasibility of recovering the mineral deposit. A preliminary survey shortly after the discovery of the deposit examines the market conditions, such as the supply and demand of the mineral, the amount of ore needed to be moved to recover a certain quantity of that mineral, and analysis of the cost associated with the operation. This pre-feasibility study determines whether the mining project is likely to be profitable; if so, a more in-depth analysis of the deposit is undertaken. After the full extent of the ore body is known and has been examined by engineers, the feasibility study examines the cost of initial capital investment, methods of extraction, the cost of operation, an estimated length of time to pay back the investment, the gross revenue and net profit margin, any possible resale price of the land, the total life of the reserve, the full value of the account, investment in future projects, and the property owner or owners' contract. In addition, environmental impact, reclamation, possible legal ramifications, and all government permitting are considered. These steps of analysis determine whether the mining company and its investors should proceed with the extraction of the minerals or whether the project should be abandoned. The mining company may decide to sell the rights to the reserve to a third party rather than develop it themselves. Alternatively, the decision to proceed with extraction may be postponed indefinitely until market conditions become favourable.
Mining operation
Mining engineers working in an established mine may work as an engineer for operations improvement, further mineral exploration, and operation capitalization by determining where in the mine to add equipment and personnel. The engineer may also work in supervision and management or as an equipment and mineral salesperson. In addition to engineering and operations, the mining engineer may work as an environmental, health, and safety manager or design engineer.
The act of mining requires different methods of extraction depending on the mineralogy, geology, and location of the resources. Characteristics such as mineral hardness, the mineral stratification, and access to that mineral will determine the method of extraction.
Generally, mining is either done from the surface or underground. Mining can also occur with surface and covert operations on the same reserve. Mining activity varies as to what method is employed to remove the mineral.
Surface mining
Surface mining comprises 90% of the world's mineral tonnage output. Also called open pit mining, surface mining removes minerals in formations near the surface. Ore retrieval is done by material removal from the land in its natural state. Surface mining often alters the land's characteristics, shape, topography, and geological makeup.
Surface mining involves quarrying and excavating minerals through cutting, cleaving, and breaking machinery. Explosives are usually used to facilitate breakage. Hard rocks such as limestone, sand, gravel, and slate are generally quarried into benches.
Using mechanical shovels, track dozers, and front-end loaders, strip mining is done on softer minerals such as clays and phosphate removed. Smoother coal seams can also be extracted this way.
With placer mining, dredge mining can also remove minerals from the bottoms of lakes, rivers, streams, and even the ocean. In addition, in-situ mining can be done from the surface using dissolving agents on the ore body and retrieving the ore via pumping. The pumped material is then set to leach for further processing. Hydraulic mining is utilized as water jets to wash away either overburden or the ore itself.
Mining process
Blasting
Explosives are used to break up a rock formation and aid in the collection of ore in a process called blasting. Blasting generally the heat and immense pressure of the detonated explosives to shatter and fracture a rock mass. The type of explosives used in mining is high explosives, which vary in composition and performance properties. The mining engineer is responsible for selecting and properly placing these explosives to maximize efficiency and safety. Blasting occurs in many phases of the mining process, such as the development of infrastructure and the production of the ore. An alternative to high explosives are Cardox blasting cartridges, invented in 1931, and extensively used from 1932 in coal mines. The cartridge contains an 'energizer' which heats liquid carbon dioxide until it ruptures a bursting disk; then, a physical explosion of the supercritical fluid.
Leaching
Leaching is the loss or extraction of certain materials from a carrier into a liquid (usually, but not always, a solvent). Mostly used in rare-earth metal extraction.
Flotation
Flotation (also spelled floatation) involves phenomena related to the relative buoyancy of minerals. It is the most widely used metal separating method.
Electrostatic separation
Separating minerals by electro-characteristic differences.
Gravity separation
Gravity separation is an industrial method of separating two components, either a suspension or dry granular mixture, where separating the components with gravity is sufficiently practical.
Magnetic separation
Magnetic separation is a process in which magnetically susceptible material is extracted from a mixture using a magnetic force.
Hydraulic separation
Hydraulic separation is a process that uses the density difference to separate minerals. Before hydraulic separation, minerals were crushed into uniform sizes; minerals with uniform sizes and densities will have different settling velocities in water, which can be used to separate target minerals.
Mining health and safety
Legal attention to health and safety in mining began in the late 19th century. In the 20th century, it progressed to a comprehensive and stringent codification of enforcement and mandatory health and safety regulation. In whatever role, a mining engineer must follow all mine safety laws.
United States
The United States Congress, through the passage of the Federal Mine Safety and Health Act of 1977, known as the Miner's Act, created the Mine Safety and Health Administration (MSHA) under the US Department of Labour. The act provides miners with rights against retaliation for reporting violations, consolidated regulation of coal mines with metallic and non-metallic mines, and created the independent Federal Mine Safety and Health Review Commission to review violations reported to MSHA.
The act codified in Code of Federal Regulations § 30 (CFR § 30) covers all miners at an active mine. When a mining engineer works at an active mine, they are subject to the same rights, violations, mandatory health and safety regulations, and compulsory training as any other worker at the mine. The mining engineer can be legally identified as a "miner".
The act establishes the rights of miners. The miner may report at any time a hazardous condition and request an inspection. The miners may elect a miners' representative to participate during an inspection, pre-inspection meeting, and post-inspection conference. The miners and miners' representatives shall be paid for their time during all inspections and investigations.
Environmental concerns
Waste and uneconomic material generated from the mineral extraction process are the primary source of pollution in the vicinity of mines. Mining activities, by their nature, cause a disturbance of the natural environment in and around which the minerals are located. Mining engineers should therefore be concerned not only with the production and processing of mineral commodities but also with the mitigation of damage to the environment both during and after mining as a result of the change in the mining area.
| Technology | Disciplines | null |
488682 | https://en.wikipedia.org/wiki/Bicycle-sharing%20system | Bicycle-sharing system | A bicycle-sharing system, bike share program, public bicycle scheme, or public bike share (PBS) scheme, is a shared transport service where bicycles are available for shared use by individuals at low cost.
The programmes themselves include both docking and dockless systems, where docking systems allow users to rent a bike from a dock, i.e., a technology-enabled bicycle rack and return at another node or dock within the system – and dockless systems, which offer a node-free system relying on smart technology. In either format, systems may incorporate smartphone web mapping to locate available bikes and docks. In July 2020, Google Maps began including bike share systems in its route recommendations.
With its antecedents in grassroots mid-1960s efforts; by 2022, approximately 3,000 cities worldwide offer bike-sharing systems, e.g., Dubai, New York, Paris, Mexico City, Montreal and Barcelona.
History
The first bike sharing projects were initiated by various sources, such as local community organizations, charitable projects intended for the disadvantaged, as way to promote bicycles as a non-polluting form of transportation – and bike-lease businesses.
The earliest well-known community bicycle program was started in the summer of 1965 by Luud Schimmelpennink in association with the group Provo in Amsterdam, the Netherlands. the group Provo painted fifty bicycles white and placed them unlocked in Amsterdam for everyone to use freely. This so-called White Bicycle Plan () provided free bicycles that were supposed to be used for one trip and then left for someone else. Within a month, most of the bikes had been stolen and the rest were found in nearby canals. The program is still active in some parts of the Netherlands, e.g., at Hoge Veluwe National Park where bikes may be used within the park. It originally existed as one in a series of White Plans proposed in the street magazine produced by the anarchist group PROVO. Years later, Schimmelpennink admitted that "the Sixties experiment never existed in the way people believe" and that "no more than about ten bikes" had been put out on the street "as a suggestion of the bigger idea." As the police had temporarily confiscated all of the White Bicycles within a day of their release to the public, the White Bicycle experiment had actually lasted less than one month.
Ernest Callenbach's novel Ecotopia (1975) illustrated the idea. In the utopian novel of a society that does not use fossil fuels, Callenbach described a bicycle sharing system which is available to inhabitants and is an integrated part of the public transportation system.
To prevent thefts, bike sharing programs gravitated to smart card control systems.' One of the first 'smart bike' programs was the Grippa™ bike storage rack system used in Portsmouth (UK)'s Bikeabout system. The Bikeabout scheme was launched in October 1995 by the University of Portsmouth, UK as part of its Green Transport Plan in an effort to cut car travel by staff and students between campus sites. Funded in part by the EU's ENTRANCE program, the Bikeabout scheme was a "smart card" fully automated system. For a small fee, users were issued magnetic striped 'smart cards' readable at a covered 'bike store' kiosk, unlocking the bike from its storage rack. Station-located CCTV cameras limited vandalism. On arrival at the destination station, the smart card unlocked cycle rack and recorded the bike's return, registering if the bike was returned with damage or if the rental time exceeded a three-hour maximum. Implemented with an original budget of approximately £200,000, the Portsmouth Bikeabout scheme was never very successful in terms of rider usage, in part due to the limited number of bike kiosks and hours of operation. Seasonal weather restrictions and concerns over unjustified charges for bike damage also imposed barriers to usage. The Bikeabout program was discontinued by the university in 1998 in favor of expanded minibus service; the total costs of the Bikeabout program were never disclosed.
One of the first community bicycle projects in the United States was started in Portland, Oregon in 1994 by civic and environmental activists Tom O'Keefe, Joe Keating and Steve Gunther. It took the approach of simply releasing a number of bicycles to the streets for unrestricted use. While Portland's Yellow Bike Project was successful in terms of publicity, it proved unsustainable due to theft and vandalism of the bicycles. The Yellow Bike Project was eventually terminated, and replaced with the Create A Commuter (CAC) program, which provides free secondhand bicycles to certain preselected low-income and disadvantaged people who need a bicycle to get to work or attend job training courses.
In 1995, a system of 300 bicycles using coins to unlock the bicycles in the style of shopping carts was introduced in Copenhagen. It was initiated by Morten Sadolin and Ole Wessung. The idea was developed by both Copenhageners after they were victims of bicycle theft one night in 1989. Copenhagen's ByCylken program was the first large-scale urban bike share program to feature specially designed bikes with parts that could not be used on other bikes. To obtain a bicycle, riders pay a refundable deposit at one of 100 special locking bike stands, and have unlimited use of the bike within a specified 'city bike zone.' The fine for not returning a bicycle or leaving the bike sharing zone exceeds US$150, and is strictly enforced by the Copenhagen police. Originally, the program's founders hoped to completely finance the program by selling advertising space on the bicycles, which was placed on the bike's frame and its solid disc-type wheels. This funding source quickly proved to be insufficient, and the city of Copenhagen took over the administration of the program, funding most of the program costs through appropriations from city revenues along with contributions from corporate donors. Since the City Bikes program is free to the user, there is no return on the capital invested by the municipality, and a considerable amount of public funds must constantly be re-invested to keep the system in service, to enforce regulations, and to replace missing bikes.
The modern wave of electronically locked bikes took off in France. In 1998 the city of Rennes France launched Velo a la cart using a magnetic card to release bicycles, which was operated by Clear Channel. Then the French advertising company, JCDecaux begain launching larger systems in Vienna (2003), Lyon (2005), and Paris (2007), among others. The Paris system captured the attention of the world and catalyzed steep growth in bikesharing systems around Europe, Asia, South America, and North America. In North America, the BIXI project (a portmanteau of the French "bicyclette" and "taxi" or "bycyle taxi") launched by the City of Montreal in 2009. It garnered a sizable ridership and the city created the Public Bike System Company to begin selling the underlying infrastructure to several other cities, including Washington D.C.'s Capital Bikeshare (2010), New York City's Citi Bike (2013), and London's "Boris bikes (2010)". The PBSC was privatised in 2014 and was later acquired by Lyft in 2022. Separately in 2018, Lyft had acquired Motivate, an operator of many BIXI-based systems. Meanwhile, the original BIXI system has been operated directly by the City of Montreal since 2014.
In 2016, the Portland Bureau of Transportation (PBOT) launched Biketown, also known as Biketown PDX, a bicycle-sharing system in Portland, Oregon. It is operated by Motivate, with Nike, Inc. as the title sponsor. At launch, the system had 100 stations and 1,000 bicycles serving the city's central and eastside neighbourhoods, with hopes to expand outward.
Bike share technology has evolved over the course of decades, and development of programs in Asia has grown exponentially. Of the world's 15 biggest public bike share programs, 13 are in China. In 2012, the biggest are in Wuhan and Hangzhou, with around 90,000 and 60,000 bikes respectively.
As of December 2016, roughly 1,000 cities worldwide have a bike-sharing program.
Categorization
Bike-sharing systems have developed and evolved with society changes and technological improvements. The systems can be grouped into five categories or generations. Many bicycle programmes paint their bicycles in a strong solid colour, such as yellow or white. Painting the bicycles helps to advertise the programme, as well as deter theft (a painted-over bicycle frame is normally less desirable to a buyer). However, theft rates in many bike-sharing programmes remain high, as most shared-use bicycles have value only as basic transport, and may be resold to unsuspecting buyers after being cleaned and repainted. In response, some large-scale bike sharing programmes have designed their own bike using specialized frame designs and other parts to prevent disassembly and resale of stolen parts.
Staffed stations
Short-term checkout
Also known as bicycle rental, bike hire or zero generation. In this system a bicycle can be rented or borrowed from a location and returned to that location. These bicycle renting systems often cater to day-trippers or tourists. This system is also used by cycling schools for potential cyclists who do not have a bicycle. The locations or stations are not automated but are run by employees or volunteers.
Regional programs have been implemented where numerous renting locations are set up at railway stations and at local businesses (usually restaurants, museums and hotels) creating a network of locations where bicycles can be borrowed from and returned (e.g. ZweiRad FreiRad with at times 50 locations). In this kind of network for example a railway station master can allocate a bicycle to a user that then returns it at a different location, for example a hotel. Some such systems require paying a fee, and some do not. Usually the user will be registered or a deposit will be left by the renting facility. The EnCicla Bike Share System in Medellín on its inception in 2011 had 6 staffed locations. It later grew to 32 automatic and 19 staffed stations making it a hybrid between a zero generation and third generation system.
Long-term checkout
Sometimes known as bike library systems, these bicycles may be lent free of charge, for a refundable deposit, or for a small fee. A bicycle is checked out to one person who will typically keep it for several months, and is encouraged or obliged to lock it between uses. A disadvantage is a lower usage frequency, around three uses per day on average as compared to 2 to 15 uses per day typically experienced with other bike-sharing schemes. Advantages of long-term use include rider familiarity with the bicycle, and constant, instant readiness.
The bicycle can be checked out like a library book, a liability waiver can be collected at check-out, and the bike can be returned any time. For each trip, a Library Bike user can choose the bike instead of a car, thus lowering car usage. The long-term rental system generally results in fewer repair costs to the scheme administrator, as riders are incentivised to obtain minor maintenance in order to keep the bike in running order during the long rental period. Most of the long-term systems implemented to date are funded solely through charitable donations of second-hand bicycles, using unpaid volunteer labour to maintain and administer the bicycle fleet. While reducing or eliminating the need for public funding, such a scheme imposes an outer limit to program expansion. The Arcata Bike Library, in California, has loaned over 4000 bicycles using this system.
White bikes
Also known as free bikes, unregulated or first generation. In this type of programme the bicycles are simply released into a city or given area for use by anyone. In some cases, such as a university campus, the bicycles are only designated for use within certain boundaries. Users are expected to leave the bike unlocked in a public area once they reach their destination. Depending on the quantity of bicycles in the system availability of such bicycles can suffer because the bikes are not required to be returned to a centralised station. Such a system can also suffer under distribution problems where many bicycles end up in a valley of a city but few are found on the hills of a city. Since parked and unlocked bikes may be taken by another user at any time, the original rider might have to find an alternative transport for the return trip. This system does away with the cost of having a person allocating a vehicle to a user and it is the system with the lowest hemmschwelle or psychological barrier for a potential user. However, bicycle sharing programs without locks, user identification, and security deposits have also historically suffered loss rates from theft and vandalism. Many initiatives have been abandoned after a few years (e.g. Portland's Yellow Bike Project was abandoned after 3 years), while others have been successful for decades (e.g. Austin's Yellow Bike Project active since 1997). Most of these systems are based around volunteer work and are supported by municipalities. Bicycle repair and maintenance are done by a volunteer project or from the municipality contracted operator but also can be, and sometimes is, completed by individual users who find a defect on a free bike.
Coin deposit stations
Also known as Bycykel or as second generation, this system was developed by Morten Sadolin and Ole Wessung of Copenhagen after both were victims of bicycle theft one night in 1989. They envisioned a freely available bicycle sharing system that would encourage spontaneous usage and also reduce bicycle theft. The bicycles, designed for intense utilitarian use with solid rubber tires and wheels with advertising plates, have a slot into which a shopping cart return key can be pushed. A coin (in most versions a 20 DKK or 2 EUR coin) needs to be pushed into the slot to unlock the bike from the station. The bicycle can thus be borrowed free of charge and for an unlimited time and the deposit coin can be retrieved by returning the bicycle to a station again. Since the deposit is a fraction of the bike's cost, and user is not registered this can be vulnerable to theft and vandalism. However, the distinct Bycykel design, well known to the public and to the law authorities does deter misuse to a degree. Implemented systems usually have a zone or area where it is allowed to drive in. The first coin deposit (small) systems were launched in 1991 in Farsø and Grenå, Denmark, and in 1993 in Nakskov, Denmark with 26 bikes and 4 stations. In 1995 the first large-scale 800 bike strong second generation bike-sharing program was launched in Copenhagen as Bycyklen. The system was further introduced in Helsinki (2000–2010) and Vienna in (2002) and in Aarhus 2003.
Automated stations
Also known as docking stations bicycle-sharing, or membership bicycles or third generation consist of bicycles that can be borrowed or rented from an automated station or "docking stations" or "docks" and can be returned at another station belonging to the same system. The docking stations are special bike racks that lock the bike, and only release it by computer control. Individuals registered with the program identify themselves with their membership card (or by a smart card, via cell phone, or other methods) at any of the hubs to check out a bicycle for a short period of time, usually three hours or less. In many schemes the first half-hour is free. In recent years, in an effort to reduce losses from theft and vandalism, many bike-sharing schemes now require a user to provide a monetary deposit or other security, or to become a paid subscriber. The individual is responsible for any damage or loss until the bike is returned to another hub and checked in.
Some cities allow to use the same card as for bus and rail transport to unlocks the bicycles.
This system was developed as Public Velo by Hellmut Slachta and Paul Brandstätter from 1990 to 1992, and first implemented in 1996 by the University of Portsmouth and Portsmouth City Council as Bikeabout with a magnetic card used by the students and on 6 June 1998 in Rennes as LE vélo STAR, a public city network with 200 bikes, 25 stations and electronic identification of the bikes or in Oslo in 2001. The smart card contactless technology was experimented in Vienna (Citybike Wien) and implemented at a large scale in 2005 in Lyon (Vélo'v) and in 2007 in Paris (Vélib'). Since then over 1000 bicycle sharing system of this generation have been launched. The countries with the most dock based systems are Spain (132), Italy (104), and China (79). , public bike share systems were available in 50 countries on five continents, including 712 cities, operating approximately 806,200 bicycles at 37,500 stations. , the Wuhan and Hangzhou Public Bicycle bike-share systems in China were the largest in the world, with around 90,000 and 60,000 bicycles respectively. By 2013, China had a combined fleet of 650,000 public bikes.
This bicycle-sharing system saves the labour costs of staffed stations (zero generation), reduces vandalism and theft compared to first and second generation systems by registering users but requires a higher investment for infrastructure compared to fourth generation dockless bikes. Third generation systems also allow adapting docking stations as recharging stations for E-bike sharing.
Dockless bikes
Also known as Call a Bike, free floating bike or fourth generation, the dockless bike hire systems consist of a bicycle with a lock that is usually integrated onto the frame and does not require a docking station. The earliest versions of this system consisted of for-rent-bicycles that were locked with combination locks and that could be unlocked by a registered user by calling the vendor to receive the combination to unlock the bicycle. The user would then call the vendor a second time to communicate where the bicycle had been parked and locked. This system was further developed by Deutsche Bahn in 1998 to incorporate a digital authentication codes (that changes) to automatically lock and unlock bikes. Deutsche Bahn launched Call a Bike in 2000, enabling users to unlock via SMS or telephone call, and more recently with an app. Recent technological and operational improvements by telephones and GPSs have paved the way for dramatic increase of this type of private app driven "dockless" bicycle-sharing system. In particular in China, Ofo and Mobike have become the world's largest bike share operators with millions of bikes spread over 100 cities. Today dockless bike shares are designed whereby a user need not return the bike to a kiosk or station; rather, the next user can find it by GPS. Over 30 private companies have started operating in China. However, the rapid growth vastly outpaced immediate demand and overwhelmed Chinese cities, where infrastructure and regulations were not prepared to handle a sudden flood of millions of shared bicycles.
Not needing docking stations that may require city planning and building permissions, the system spread rapidly on a global scale. At times dockless bike-sharing systems have been criticized as rogue systems instituted without respect for local authorities. In many cities entrepreneurial companies have independently introduced this system, despite a lack of adequate parking facilities. City officials lack regulation experience for this mode of transportation and social habits have not developed either. In some jurisdictions, authorities have confiscated "rogue" dockless bicycles that are improperly parked for potentially blocking pedestrian traffic on sidewalks and in other cases new laws have been introduced to regulate the shared bikes.
In some cities Deutsche Bahn's Call a Bike has Call a Bike fix system, which has fixed docking stations versus the flex dockless version, some systems are combined into a hybrid of third and fourth generation systems. Some Nextbike systems are also a 3rd and 4th generation hybrid. With the arrival of dockless bike shares, there were in 2017 over 70 private dockless bikeshares operating a combined fleet of 16 million share bikes according to estimates of Ministry of Transport of China. Beijing alone has 2.35 million share bikes from 15 companies.
In the United States, many major metropolitan areas are experimenting with dockless bikeshare systems, which have been popular with commuters but subject to complaints about illegal parking.
Goals
People use bike-share for various reasons. Cost and time are primary motivators for using bike-sharing programs, in particular the perceived cost of travel and time saved traveling. Some who would otherwise use their own bicycle have concerns about theft, vandalism, parking, storage, and maintenance.
Sustainable alternative for short trips
Most large-scale urban bike sharing programmes have numerous bike check-out stations, and operate much like public transit systems, catering to tourists and visitors as well as local residents. Their central concept is to provide free or affordable access to bicycles for short-distance trips in an urban area as an alternative to private vehicles, thereby reducing congestion, noise, and air pollution. According to research in 2016, the bike sharing system in Shanghai saved 8,358 tonnes of petrol and decreased carbon dioxide and NOx emissions by 25,240 and 64 tonnes, respectively. The research also stated that bike sharing system has great potential to reduce energy consumption and emissions based on its rapid development.
Last mile problem
Bicycle-sharing systems have also been cited as a way to solve the "last mile" problem of public transit networks. According to a research conducted on YouBike system in Taipei, on 2014, the bike sharing system in residential area are more popular, and as a first/last mile of transport mode to and from the station to their desired locations. However, dock systems, serving only stations, resemble public transit and have therefore been criticized as less convenient than a privately owned bicycle used door-to-door.
Operation
Bicycle-sharing systems are an economic good, and are generally classified as a private good due to their excludable and rivalrous nature. While some bicycle-sharing systems are free, most require some user fee or subscription, thus excluding the good to paying consumers. Bicycle-sharing systems also provide a discrete and limited number of bikes, whose distribution can vary throughout a city. One person's usage of the good diminishes the ability of others to use the same good. Nonetheless, the hope of many cities is to partner with bike-share companies to provide something close to a public good. Public good status may be achieved if the service is free to consumers and there are a sufficient number of bicycles such that one person's usage does not encroach upon another's use of the good.
Partnership with public transport sector
In a national-level programme that combines a typical rental system with several of the above system types, a passenger railway operator or infrastructure manager partners with a national cycling organisation and others to create a system closely connected with public transport. These programmes usually allow for a longer rental time of up to 24 or 48 hours, as well as tourists and round trips. In some German cities the national rail company offers a bike rental service called Call a Bike.
In Guangzhou, China, the privately operated Guangzhou Bus Rapid Transit system includes cycle lanes, and a public bicycle system.
In some cases, like Santander Cycles in London, the bicycle sharing system is owned by the public transport authority itself.
In other cases, like Youbike in Taipei, Taiwan, the bicycle sharing system is built by a private company partner with the public transport sector through BOT mode. To be more specific in this case, it is offered by the Taipei City Department of Transportation in a BOT collaboration with local manufacturer Giant Bicycles.
Partnership with other modes of public transport
In many cities over the world, bike sharing system is connected to other public transportation. It is usually hoped to complement the shortcomings in the greater public transport system. Sometimes, in order to encourage residents to use public transport system, local government will give discount on transferring between bike sharing system and other public transports.
Medellin
The city of Medellin is home to 3.4 million inhabitants in 173 km2 and has long faced infrastructural mobility challenges. EnCicla is a bike sharing system in the city of Medellin (Colombia, South America). The bike sharing system is connected to other modes of transportation, such as the Metro.
In 2010, three EAFIT students (Lina Marcela López, José Agusto Ocampo, and Felipe Gutiérrez) developed the idea of the EnCicla bike sharing system as part of their final project. The implementation of the system was decided in operation in August 2012, with the subsequent pilot program confirming its prospects for success. EAFIT advocated for the city to lead the system. This was implemented accordingly, resulting in the inclusion of EnCicla in the agenda of the city of Medellin and its incorporation into the transportation network. In this regard, EnCicla consists of a mixture of shared, as well as separated, bike lanes on the roadway. In the first 3 months after the official launch, 15,700 bicycle rentals took place, with usage picking up sharply in subsequent months and years. In Medellin, an attempt was made to solve the demand problem with statistical analysis using historical data. The result of this analysis was the establishment of a heterogeneous bicycle fleet, with a minimum and maximum number for each station.
In total, in Medellin there exist more than 90 stations in 7 zones, with 13 connected to other transport systems. Since inception, more than 13 million bicycles have been rented by the approximately 9,100 active members. In this context, the most frequently used stations are located in the western zone, near universities and colleges. These stations are located near train stations, which means that there is a high volume of people. To use EnCicla, citizens must register on the official website. In general, the system can be used free of charge by anyone 16 years of age or older and is available from 5:30-22:00 during the week and from 6:30-21:00 on Saturdays. Local residents must register through EnCicla's website prior to use, and tourists have the option of renting a bicycle using their passport.
The establishment of EnCicla in recent years has helped relieve the complex transportation system in Medelin. However, the repositioning of bicycles at stations results in increased CO2 emissions, which run counter to the environmental importance of the project. In parallel, various activities have been carried out to promote the establishment of the system. These include a program that gives people over 8 years of age the opportunity to improve their knowledge and skills in cycling.
Taipei Metropolitan Area
YouBike, a bike sharing system in Taipei–Keelung metropolitan area, Taiwan, has automated stations near all Taipei Metro stations. The integration of YouBike stations and Taipei Metro aims at solving the "last mile" problem, thus improving transit accessibility and usability. It is hoped that YouBike could complement the shortcomings in the greater public transport. Commuters can check in or check out YouBikes near the metro stations to catch connections from the station to the destination.
Transfer Discount Offered for Commuters
Starting 30 March 2021, passengers renting a YouBike from any YouBike station in the Taipei–Keelung metropolitan area receive a discount of NT$5 when using their EasyCard to transfer between YouBike and Taipei Metro, local buses (except buses that charge by distance) or Danhai LRT within one hour. Plus, the trip is only eligible for a discount when the transfer is direct. Commuters shall not utilise other means of transportation, such as Taiwan Railways, Maokong Gondola, long-distance buses, Taiwan High Speed Rail, Taoyuan Metro, or taxis.
Transfer Behavior
According to the analysis of YouBike rental and its Taipei MRT (Taipei Rapid Transit System) transfer behavior from the Department of Transportation, New Taipei City Government, YouBike has already become an important feeder mode for metro commuters: up to 55% of the subjects (the commuters who ever utilise YouBike during September, 2015) transfer by YouBike before or after taking the Metro. Adopting the YouBike and MRT transaction data of EasyCard in New Taipei City in November, 2016, almost all popular YouBike stations can be found next to the Taipei metro stations. Furthermore, transfer analysis depending on the YouBike and MRT data indicates that, the transfer ratio of loyal users (who utilise YouBike more than five times per week) is up to 60%.
Seoul Metropolitan Area
Sharing bicycles in South Korea are called 'Ddareungi' in Seoul capital area. Ddareungi is a sharing bicycle operated throughout Seoul. It is an unmanned sharing bicycle rental service that started pilot operation in 2014 and officially operated in October 2015.
The 1-hour pass for Ddareungi is KRW 1000(Approximate 1 USD), and to prevent theft, an additional charge of KRW 1000 per 30 minutes is charged for exceeding the usage time.
Transit Mileage
Transit Mileage is a benefit that can only be received by 365-day commuter pass users. If someone uses public transportation within 30 minutes of returning the bicycle, the mileage is accumulated. If it is difficult to travel by bus or subway, the section can be replaced with Ddareungi.
Bicycle Driving Ability Certification System Fee Benefits
Bicycle driving ability certification system requires completion of bicycle safety education, if a person passes both the written and practical exams, that person will receive certification and part of the Ddareungi usage fee can be reduced for two years.
QR Code Lock
From 1 March 2020, QR Code Lock was introduced as a method of renting and returning by recognizing QR codes. It is convenient because it can be rented or returned with a single scan by using a QR code-type locking device. When renting a bicycle, purchase a voucher from the bicycle app and scan the QR code on the bicycle to rent, and the lock is automatically unlocked and can be used immediately. It can return and rent a bicycle anywhere without going to a bicycle rental booth.
Sprout Ddareungyi
Existing sharing bicycles can only be used by those over the age of 15, so Sprout Ddareungyi, which can be rented from the age of 13 and older, has been launched in Seoul. The government released a policy for public bicycles with reduced size and weight compared to the existing Ddareungi bicycles so that even small-sized people, such as teenagers and the elderly, could use them conveniently.
The number of users of Seoul's public bicycle 'Ddareungi' has exceeded 3 million. It is used by about one in three Seoul citizens. General citizens have a high rate of use during commuting hours on weekdays, except on weekends, so after using public transportation such as Seoul Metropolitan Subway as well as Seoul Buses, when it is an ambiguous distance to use public transportation anymore, citizens use public bicycles near subway stations to move the most. In particular, considering that rentals and returns are made at rental stations near subway stations, citizens frequently use Ddareungi that are deployed in subway stations. To analyse, if the number of cases is classified based on the number of Ddareungi rental stations near subway stations in 2021, exit 1 of the Ttukseom Park area of Hangang Park, which is the most used in Seoul, is the first with 602 rentals, After that, Express Bus Terminal Station and Lotte World Tower's Jamsil Station Exit 2 followed. It is analysed that the most frequently used Ttukseom area, exit 1, is usually used by citizens who enjoy leisure at Hangang Park except during rush hour.
Hamburg
The bicycle sharing system "StadtRAD" of Hamburg (Germany) was launched in 2009 and now includes 3,100 bicycles and 20 cargo bikes. The infrastructure includes 250 fixed stations distributed throughout the city. With the help of the app, it is possible to rent up to two bikes. For an annual fee of €5, the first 30 minutes of each rental are free. A total of 500,000 people are registered with the app.
"StadtRAD" is an integral part of the city's mobility transition. The aim here is to reduce motorized private transport by strengthening Public transport, making it easier to switch between different modes of transport and developing the city into a bicycle city. To achieve this, the share of cyclists in total traffic should increase to 25%.
The administrative responsibility for implementing the mobility strategy is assumed by the Alliance for Cycling, which is assigned to the Authority for Transport and Mobility Change. Its task is to make the transport infrastructure bicycle-friendly by promoting the construction of bike and ride facilities, making subway stops barrier-friendly and expanding bicycle routes. In addition, it should be made easier to take bicycles on buses and trains, and traffic safety should be strengthened through traffic education in schools. Another focus is on the interconnection of the different mobility offers through the Switchh app. Citizens can switch between Carsharing, Motorized scooter or offers from the HVV by making reservations or bookings. The integration of "StadtRAD" is planned to take place in 2022.
For the practical implementation of the bicycle sharing system Stadtrad, the city has contracted "Deutsche Bahn Connect". Deutsche Bahn Connect is committed to setting up and operating a public bike sharing system with fixed rental stations within the city boundaries.
In a study conducted by the University of Hamburg, users state that they use the Stadtrad mainly for leisure (55.9%) and regularly spend time in the city center (89.9%). The frequency of use is several times a month (24.9%) and several times a half year (24.9%). In addition, the study shows that the city bike has a positive image among users, they are largely satisfied with the service and recognize environmental and health benefits. Another study by the authors shows that implementing a green public service increases both perceived social and environmental value. Perceived social and environmental values have a positive influence on the user's green attitude and intentions.
At the same time, however, the need for a reporting system for sharing systems is emphasized, which ideally should be standardized and comparable with other regions. Especially for outsourced projects, monitoring and control processes must be implemented to ensure consistent quality. This includes the use of advanced prediction models to balance the load effectively by forecasting bike-sharing activity and addressing underflow and overflow issues dynamically. In addition to the environmental benefits, financial and time constraints must also be considered in large urban planning projects. Hamburg has had increasing spending on bicycle infrastructure since 2011 and spent 15 million euros on it in 2017.
Partnership with car park operators
Some car park operators such as Vinci Park in France lend bikes to their customers who park a car.
Partnership with car-share operations
City CarShare, a San Francisco-based non-profit, received a federal grant in 2012 to integrate electric bicycles within its existing car-sharing fleet. The program is set to launch before the end of 2012 with 45 bikes.
Financing
The financing of bicycle-sharing system have been maintained by a combination of fees, volunteer, charity, advertisements, business interest groups and government subsidies. The international expansion dockless bicycles in mid 2010s has been financed by investment capital.
User fees
User rent fees may range from the equivalent of US$0.50 to 30.00 per day, rent fees for 15- or 20-minute intervals can range from a few cents to 1.00. Many bike-share systems offer subscriptions that make the first 30–45 minutes of use either free or very inexpensive, encouraging use as transportation. This allows each bike to serve several users per day but reduces revenue. Monthly or yearly membership subscriptions and initial registration fees may apply. To reduce losses from theft often users are required to commit to temporary deposit via a credit card or debit card. If the bike is not returned within the subscription period, or returned with significant damage, the bike sharing operator keeps the deposit or withdraws money from the user's credit card account. operated by private companies as is the case in most cities in China.
US
New York rental rates are among the highest in the world, as of the Citi Bike program's launch in July 2012, but this is of course subject to change. The cost of annual membership in the US varies between $100 and about $220.
Europe
Bike riders shared in Europe usually pay between €0.50 to €1 per trip, and an average of €10–12 for a full day cycling.
Volunteer work
Many first and second generation bicycle sharing programs were and are community run organisations as "Community Bike programmes", as done in IIT Bombay. Often maintenance and repair is performed by unpaid volunteers that complete this work in their own free time.
Charity sources
Charity fundraising drives and charitable organisations have and do support bicycle sharing programs, including Rotary Clubs and Lions Clubs.
Advertisement revenue
Second and third generation schemes in the 90s already prominently included advertising opportunities on the individual bikes in form of advertisement areas on the wheels or frame. Other schemes are completely branded according to a sponsor, notable example London's bike share which was originally branded and sponsored by Barclays Bank and subsequently by Santander UK.
Several European cities, including the French cities of Lyon and Paris as well as London, Barcelona, Stockholm and Oslo, have signed contracts with private advertising agencies (JCDecaux in Brussels, Lyon, Paris, Seville, Dublin and Oslo; Clear Channel in Stockholm, Barcelona, Antwerp, Perpignan and Zaragoza) which supply the city with thousands of bicycles free of charge (or for a minor fee). In return, the agencies are allowed to advertise both on the bikes themselves and in other select locations in the city, typically in the form of advertising on stations or the bicycles themselves.
Government subsidies
Municipalities have operated and do operate bicycle share systems as a public service, paying for the initial investment, maintenance and operations if it is not covered by other revenue sources. Governments can also support bicycle share programs in forms of one time grants (often to buy a set of bicycles), yearly of monthly subsidies, or by paying part of the employee wages (example in repair workshops that employee long-term unemployed persons). Many of the membership-based systems are operated through public-private partnerships. Some schemes may be financed as a part of the public transportation system (for example Smoove). In Melbourne the government subsidises the sale of bicycle helmets to enable spontaneous cyclists comply with the mandatory helmet laws.
Harvesting of user-data
GPS traceable vehicle commute patterns and usage habits present valuable data for government agencies, marketing companies or researchers. Strong commuter patterns can be filtered out and potential transportation services (e.g. commuter bus) can be tailored to existing demand. Potential audiences can be better assessed and understood.
Usage patterns
Most bike-sharing systems allow the bicycles to be returned to any station in the system, which facilitates one-way trips because the users do not need to return the bicycles to the origin. Thus, one bike may take 10–15 rides a day with different users and can be ridden up to a year (as in Vélo'v in Lyon, France). Each bike has at least one rides with one unique user per day which indicates that in 2014 there were a minimum of at least 294 million unique bike share cyclists worldwide (806,200 bicycles × 365) although some estimates are much.
It was found—in cities like Paris and Copenhagen—that to have a major impact there had to be a high density of available bikes. Copenhagen has 2500 bikes which cannot be used outside the zone of the city centre (a fine of DKK 1000 applies to any user taking bikes across the canal bridges around the periphery). Since Paris's Vélib' programme operates with an increasing fee past the free first half-hour, users have a strong disincentive to take the bicycles out of the city centre. The distance between stations is only in inner city areas.
in US, male users of bike sharing made up for more than 80% of total trips made in 2017.
A study published in 2015 in the journal Transportation concludes that bike sharing systems can be grouped into behaviourally similar categories based upon their size. Cluster analysis shows that larger systems have different usage patterns in different stations, whilst in smaller systems the different stations have similar daily utilization patterns.
Global distribution of bike-sharing systems
Economic impact
Bike-share programs generate a number of economic externalities, both positive and negative. The positive externalities include reduction of traffic congestion and pollution, while the negative externalities can include degradation of urban aesthetic environment and reduction of parking. Furthermore, bike-share programs have pecuniary effects. Some of these economic externalities (e.g. reduced congestion) can be systematically evaluated using empirical data, and therefore may be internalized through government subsidy. On the other hand, "nuisance" externalities (e.g. street and sidewalk clutter) are more subjective and harder to quantify, and may not be able to be internalized.
Positive effects
Less traffic congestion
A primary goal of bicycle-sharing systems has been to reduce traffic congestion, particularly in large urban areas. Some empirical evidence indicates that this goal has been achieved to varying degrees in different cities. A 2015 article in Transport Reviews examined bike-share systems in five cities, including Washington, D.C., and Minneapolis. The article found that in D.C., individuals substituted bike-share rides for automobile trips 8 percent of the time, and almost 20 percent of the time in Minneapolis. A separate study on Washington, D.C.'s Capital Bikeshare found that the bike-share program contributed a 2 to 3 percent reduction in traffic congestion within the evaluated neighborhood. 2017 studies in Beijing and Shanghai have linked the massive increase of dockless bike shares to the decrease in the number of private automobile trips that are less than five kilometres. In Guangzhou, the arrival of dockless bike shares had a positive impact in the growth of cycling mode share.
Less pollution
Not only do bike-share systems intend to reduce traffic congestion, they also aim to reduce air pollution through decreased automobile usage, and indirectly through the reduction of congestion. The study on D.C.'s Capital Bikeshare estimated that the reduction in traffic congestion would be equivalent to roughly $1.28 million in annual benefits, accrued through the reduction in congestion-induced CO2 emissions. A separate study of transportation in Australia estimated that 1.5 kilograms of CO2 equivalent emissions are avoided by an urban resident who travels 5 kilometers by cycling rather than by car during rush hour periods.
Healthy transport
Bicycle-sharing systems have been shown to have a strong net positive health effect. Cycling is a good way for exercise and stress relief. It can increase recreation and improve sociability of a city, which make people live more happy and relaxed. The report from Centers for Disease Control and Prevention (CDC) point out that cycling also help preventing disease like obesity, heart disease (can reduce up to 82%) and diabetes (can reduce up to 58%). Therefore, bicycle-sharing systems has a positive effect on mental and physical health, which attract more people to use. (Demand increase)
Reduced car parking
Bike-share programs, especially the earlier services that required docking areas along urban streets, may encroach upon the space available for on-street car parking and other auto-centric uses. While some argue that this is a negative, it is generally considered a positive side effect, since it helps further the transition away from car-dependency.
Negative effects
Urban clutter
In some cities, the many dockless bike-share bicycles have cluttered streets and sidewalks, degrading the urban aesthetic environment and blocking pedestrian traffic. In particular, cycles on Chinese city streets have created sections of clogged sidewalks no longer walkable, and piles of illegally parked bicycles.
Due to the vehicles being left in the public right of way, or abandoned obstructing pedestrians, the dockless vehicles have been called "litter bikes".
Dockless cycles left randomly on public footpaths may impede access for wheelchair users and others who use mobility aids, and may be dangerous to people with visual impairments.
Pecuniary effects
As bicycle-sharing systems continue to grow and provide an affordable alternative for commuters, the relatively low price of these services may induce competitors to offer lower prices. For instance, municipal public transit organizations may lower prices for buses or subways to continue to compete with bike-share systems. Pecuniary effects may even extend to bicycle manufacturers and retailers, where these producers might reduce prices of bicycles and other complementary goods (e.g. helmets, lights). However, empirical research is needed to test these hypotheses.
Internalization of externalities
Public-private partnerships
In public economics, there is a role for government intervention in a market if market failures exist, or in the case of redistribution. As several studies have found, bike-share programs appear to produce net positive externalities in reduced traffic congestion and pollution, for example. The bike-sharing market does not produce at the social optimum, justifying the need for government intervention in the form of a subsidy for the provision of this good in order to internalize the positive externality. Many cities have adopted public-private partnerships to provide bike-shares, such as in Washington, D.C., with Capital Bikeshares. These partially government-funded programs may serve to better provide the good of bike-shares.
Dangers of over-supply
Many bike-share companies and public-private partnerships aim to supply shared bicycles as a public good. In order for bike-shares to be a public good, they must be both non-excludable and non-rival. Numerous bike-share programs already offer their services partly for free or at least at very low prices, therefore nearing the non-excludable requirement. However, in order to achieve the non-rival requirement, shared bicycles must be supplied at a certain density within an urban area. There are numerous challenges with attaining non-rivalry, for instance, redistribution of bicycles from low-demand regions to regions with high-demand. Mobike, a China-based company, has addressed this problem by paying their users to ride their bikes from low-demand areas to high-demand areas. Citi Bike in New York City has a similar "Bike Angel" program to give discounts and prizes to balancers.
Other companies such as oBike have introduced a points system to penalize negative behavior, namely, illegal parking of shared bicycles. Economists speculate that a combination of efficient pricing with well-designed regulatory policies could significantly mitigate problems of over-supply and clutter.
The Chinese bicycle-sharing market demonstrated the danger of oversupply in 2018. Companies took advantage of unclear regulations in the preceding years to introduce millions of shared bikes to the country's cities. Users were not educated in how to use the systems properly and in many cases treated them as disposable, parking them anywhere. City governments were forced to impound the abandoned bikes when they blocked public thoroughfares, and millions of bikes went directly to junkyards after the companies that owned them went bankrupt.
Health impacts
A study published in the American Journal of Public Health reports observing an increase in cycling and health benefits where bicycle sharing systems are run. In the United States, bikesharing programs have proliferated in recent years, but collision and injury rates for bikesharing are lower than previously computed rates for personal bicycling; at least two people have been killed while using a bike share scheme.
There is also considerable evidence that bike-share programs must be adopted in tandem with city infrastructure, namely, the creation of bike lanes. A 2012 study published in the American Journal of Public Health found that Toronto's cyclists were 30–50% more likely to be involved in an accident on major roads without cycle lanes than on those with.
Criticism
Despite their theoretical and observed benefits, bike-share programs have come under attack as their presence has grown throughout the world. Much of this criticism has focused on the use of public funding – concerned critics posit that the use of tax money for bike-share programs should instead be diverted towards other services that more residents use on a daily basis. However, this argument relies on a faulty assumption that taxpayer money is a significant source of bike-share funding. An analysis by People for Bikes, an organization that advocates for new and safe bike infrastructure, found that public investment in Salt Lake City's Greenbike and Denver's B-Cycle programs was significantly less than traditional public transit (e.g. bus or rail) in those same cities, on a per-trip basis. Both Greenbike and B-Cycle's publicly funded subsidies amount to 10 percent or less of the total cost of one trip. In contrast, Salt Lake City's bus and rail system (UTA) relies on 80 percent public funding for a single trip.
Other critics claim that bike-share programs fail to reach more low-income communities. Some efforts have attempted to address this issue, such as New York City's Citi Bike's discounted membership program, which is aimed at increasing ridership among low-income residents. However, around 80 percent of study respondents reported that they had no knowledge of the program's discount.
A further criticism describes increasing discriminatory technical and organizational hurdles. In addition to registration—providing addresses—or security deposits—of money or bank card data—many systems require smartphones with certain operating systems and user accounts, usually by Apple or Google, or even a permanent or temporary mobile data connection for unlocking and returning the bicycles. Others offer the same functions via SMS, telephone, or a previously purchased chip card.
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488743 | https://en.wikipedia.org/wiki/Carbon%20monoxide%20poisoning | Carbon monoxide poisoning | Carbon monoxide poisoning typically occurs from breathing in carbon monoxide (CO) at excessive levels. Symptoms are often described as "flu-like" and commonly include headache, dizziness, weakness, vomiting, chest pain, and confusion. Large exposures can result in loss of consciousness, arrhythmias, seizures, or death. The classically described "cherry red skin" rarely occurs. Long-term complications may include chronic fatigue, trouble with memory, and movement problems.
CO is a colorless and odorless gas which is initially non-irritating. It is produced during incomplete burning of organic matter. This can occur from motor vehicles, heaters, or cooking equipment that run on carbon-based fuels. Carbon monoxide primarily causes adverse effects by combining with hemoglobin to form carboxyhemoglobin (symbol COHb or HbCO) preventing the blood from carrying oxygen and expelling carbon dioxide as carbaminohemoglobin. Additionally, many other hemoproteins such as myoglobin, Cytochrome P450, and mitochondrial cytochrome oxidase are affected, along with other metallic and non-metallic cellular targets.
Diagnosis is typically based on a HbCO level of more than 3% among nonsmokers and more than 10% among smokers. The biological threshold for carboxyhemoglobin tolerance is typically accepted to be 15% COHb, meaning toxicity is consistently observed at levels in excess of this concentration. The FDA has previously set a threshold of 14% COHb in certain clinical trials evaluating the therapeutic potential of carbon monoxide. In general, 30% COHb is considered severe carbon monoxide poisoning. The highest reported non-fatal carboxyhemoglobin level was 73% COHb.
Efforts to prevent poisoning include carbon monoxide detectors, proper venting of gas appliances, keeping chimneys clean, and keeping exhaust systems of vehicles in good repair. Treatment of poisoning generally consists of giving 100% oxygen along with supportive care. This procedure is often carried out until symptoms are absent and the HbCO level is less than 3%/10%.
Carbon monoxide poisoning is relatively common, resulting in more than 20,000 emergency room visits a year in the United States. It is the most common type of fatal poisoning in many countries. In the United States, non-fire related cases result in more than 400 deaths a year. Poisonings occur more often in the winter, particularly from the use of portable generators during power outages. The toxic effects of CO have been known since ancient history. The discovery that hemoglobin is affected by CO emerged with an investigation by James Watt and Thomas Beddoes into the therapeutic potential of hydrocarbonate in 1793, and later confirmed by Claude Bernard between 1846 and 1857.
Background
Carbon monoxide is not toxic to all forms of life, and the toxicity is a classical dose-dependent example of hormesis. Small amounts of carbon monoxide are naturally produced through many enzymatic and non-enzymatic reactions across phylogenetic kingdoms where it can serve as an important neurotransmitter (subcategorized as a gasotransmitter) and a potential therapeutic agent. In the case of prokaryotes, some bacteria produce, consume and respond to carbon monoxide whereas certain other microbes are susceptible to its toxicity. Currently, there are no known adverse effects on photosynthesizing plants.
The harmful effects of carbon monoxide are generally considered to be due to tightly binding with the prosthetic heme moiety of hemoproteins that results in interference with cellular operations, for example: carbon monoxide binds with hemoglobin to form carboxyhemoglobin which affects gas exchange and cellular respiration. Inhaling excessive concentrations of the gas can lead to hypoxic injury, nervous system damage, and even death.
As pioneered by Esther Killick, different species and different people across diverse demographics may have different carbon monoxide tolerance levels. The carbon monoxide tolerance level for any person is altered by several factors, including genetics (hemoglobin mutations), behavior such as activity level, rate of ventilation, a pre-existing cerebral or cardiovascular disease, cardiac output, anemia, sickle cell disease and other hematological disorders, geography and barometric pressure, and metabolic rate.
Physiology
Carbon monoxide is produced naturally by many physiologically relevant enzymatic and non-enzymatic reactions best exemplified by heme oxygenase catalyzing the biotransformation of heme (an iron protoporphyrin) into biliverdin and eventually bilirubin. Aside from physiological signaling, most carbon monoxide is stored as carboxyhemoglobin at non-toxic levels below 3% HbCO.
Therapeutics
Small amounts of CO are beneficial and enzymes exist that produce it at times of oxidative stress. A variety of drugs are being developed to introduce small amounts of CO, these drugs are commonly called carbon monoxide-releasing molecules. Historically, the therapeutic potential of factitious airs, notably carbon monoxide as hydrocarbonate, was investigated by Thomas Beddoes, James Watt, Tiberius Cavallo, James Lind, Humphry Davy, and others in many labs such as the Pneumatic Institution.
Signs and symptoms
On average, exposures at 100 ppm or greater is dangerous to human health. The WHO recommended levels of indoor CO exposure in 24 hours is 4 mg/m3. Acute exposure should not exceed 10 mg/m3 in 8 hours, 35 mg/m3 in one hour and 100 mg/m3 in 15 minutes.
Acute poisoning
The main manifestations of carbon monoxide poisoning develop in the organ systems most dependent on oxygen use, the central nervous system and the heart. The initial symptoms of acute carbon monoxide poisoning include headache, nausea, malaise, and fatigue. These symptoms are often mistaken for a virus such as influenza or other illnesses such as food poisoning or gastroenteritis. Headache is the most common symptom of acute carbon monoxide poisoning; it is often described as dull, frontal, and continuous. Increasing exposure produces cardiac abnormalities including fast heart rate, low blood pressure, and cardiac arrhythmia; central nervous system symptoms include delirium, hallucinations, dizziness, unsteady gait, confusion, seizures, central nervous system depression, unconsciousness, respiratory arrest, and death. Less common symptoms of acute carbon monoxide poisoning include myocardial ischemia, atrial fibrillation, pneumonia, pulmonary edema, high blood sugar, lactic acidosis, muscle necrosis, acute kidney failure, skin lesions, and visual and auditory problems. Carbon monoxide exposure may lead to a significantly shorter life span due to heart damage.
One of the major concerns following acute carbon monoxide poisoning is the severe delayed neurological manifestations that may occur. Problems may include difficulty with higher intellectual functions, short-term memory loss, dementia, amnesia, psychosis, irritability, a strange gait, speech disturbances, Parkinson's disease-like syndromes, cortical blindness, and a depressed mood. Depression may occur in those who did not have pre-existing depression. These delayed neurological sequelae may occur in up to 50% of poisoned people after 2 to 40 days. It is difficult to predict who will develop delayed sequelae; however, advanced age, loss of consciousness while poisoned, and initial neurological abnormalities may increase the chance of developing delayed symptoms.
Chronic poisoning
Chronic exposure to relatively low levels of carbon monoxide may cause persistent headaches, lightheadedness, depression, confusion, memory loss, nausea, hearing disorders and vomiting. It is unknown whether low-level chronic exposure may cause permanent neurological damage. Typically, upon removal from exposure to carbon monoxide, symptoms usually resolve themselves, unless there has been an episode of severe acute poisoning. However, one case noted permanent memory loss and learning problems after a three-year exposure to relatively low levels of carbon monoxide from a faulty furnace.
Chronic exposure may worsen cardiovascular symptoms in some people. Chronic carbon monoxide exposure might increase the risk of developing atherosclerosis. Long-term exposures to carbon monoxide present the greatest risk to persons with coronary heart disease and in females who are pregnant.
In experimental animals, carbon monoxide appears to worsen noise-induced hearing loss at noise exposure conditions that would have limited effects on hearing otherwise. In humans, hearing loss has been reported following carbon monoxide poisoning. Unlike the findings in animal studies, noise exposure was not a necessary factor for the auditory problems to occur.
Fatal poisoning
One classic sign of carbon monoxide poisoning is more often seen in the dead rather than the living – people have been described as looking red-cheeked and healthy. However, since this "cherry-red" appearance is more common in the dead, it is not considered a useful diagnostic sign in clinical medicine. In autopsy examinations, the appearance of carbon monoxide poisoning is notable because unembalmed dead persons are normally bluish and pale, whereas dead carbon-monoxide poisoned people may appear unusually lifelike in coloration. The colorant effect of carbon monoxide in such postmortem circumstances is thus analogous to its use as a red colorant in the commercial meat-packing industry.
Epidemiology
The true number of cases of carbon monoxide poisoning is unknown, since many non-lethal exposures go undetected. From the available data, carbon monoxide poisoning is the most common cause of injury and death due to poisoning worldwide. Poisoning is typically more common during the winter months. This is due to increased domestic use of gas furnaces, gas or kerosene space heaters, and kitchen stoves during the winter months, which if faulty and/or used without adequate ventilation, may produce excessive carbon monoxide. Carbon monoxide detection and poisoning also increases during power outages, when electric heating and cooking appliances become inoperative and residents may temporarily resort to fuel-burning space heaters, stoves, and grills (some of which are safe only for outdoor use but nonetheless are errantly burned indoors).
It has been estimated that more than 40,000 people per year seek medical attention for carbon monoxide poisoning in the United States. 95% of carbon monoxide poisoning deaths in Australia are due to gas space heaters. In many industrialized countries, carbon monoxide is the cause of more than 50% of fatal poisonings. In the United States, approximately 200 people die each year from carbon monoxide poisoning associated with home fuel-burning heating equipment. Carbon monoxide poisoning contributes to the approximately 5,613 smoke inhalation deaths each year in the United States. The CDC reports, "Each year, more than 500 Americans die from unintentional carbon monoxide poisoning, and more than 2,000 commit suicide by intentionally poisoning themselves." For the 10-year period from 1979 to 1988, 56,133 deaths from carbon monoxide poisoning occurred in the United States, with 25,889 of those being suicides, leaving 30,244 unintentional deaths. A report from New Zealand showed that 206 people died from carbon monoxide poisoning in the years of 2001 and 2002. In total carbon monoxide poisoning was responsible for 43.9% of deaths by poisoning in that country. In South Korea, 1,950 people had been poisoned by carbon monoxide with 254 deaths from 2001 through 2003. A report from Jerusalem showed 3.53 per 100,000 people were poisoned annually from 2001 through 2006. In Hubei, China, 218 deaths from poisoning were reported over a 10-year period with 16.5% being from carbon monoxide exposure.
Causes
Carbon monoxide is a product of combustion of organic matter under conditions of restricted oxygen supply, which prevents complete oxidation to carbon dioxide (CO2). Sources of carbon monoxide include cigarette smoke, house fires, faulty furnaces, heaters, wood-burning stoves, internal combustion vehicle exhaust, electrical generators, propane-fueled equipment such as portable stoves, and gasoline-powered tools such as leaf blowers, lawn mowers, high-pressure washers, concrete cutting saws, power trowels, and welders. Exposure typically occurs when equipment is used in buildings or semi-enclosed spaces.
Riding in the back of pickup trucks has led to poisoning in children. Idling automobiles with the exhaust pipe blocked by snow has led to the poisoning of car occupants. Any perforation between the exhaust manifold and shroud can result in exhaust gases reaching the cabin. Generators and propulsion engines on boats, notably houseboats, have resulted in fatal carbon monoxide exposures.
Poisoning may also occur following the use of a self-contained underwater breathing apparatus (SCUBA) due to faulty diving air compressors.
In caves carbon monoxide can build up in enclosed chambers due to the presence of decomposing organic matter. In coal mines incomplete combustion may occur during explosions resulting in the production of afterdamp. The gas is up to 3% CO and may be fatal after just a single breath. Following an explosion in a colliery, adjacent interconnected mines may become dangerous due to the afterdamp leaking from mine to mine. Such an incident followed the Trimdon Grange explosion which killed men in the Kelloe mine.
Another source of poisoning is exposure to the organic solvent dichloromethane, also known as methylene chloride, found in some paint strippers, as the metabolism of dichloromethane produces carbon monoxide. In November 2019, an EPA ban on dichloromethane in paint strippers for consumer use took effect in the United States.
Prevention
Detectors
Prevention remains a vital public health issue, requiring public education on the safe operation of appliances, heaters, fireplaces, and internal-combustion engines, as well as increased emphasis on the installation of carbon monoxide detectors. Carbon monoxide is tasteless, odourless, and colourless, and therefore can not be detected by visual cues or smell.
The United States Consumer Product Safety Commission has stated, "carbon monoxide detectors are as important to home safety as smoke detectors are," and recommends each home have at least one carbon monoxide detector, and preferably one on each level of the building. These devices, which are relatively inexpensive and widely available, are either battery- or AC-powered, with or without battery backup. In buildings, carbon monoxide detectors are usually installed around heaters and other equipment. If a relatively high level of carbon monoxide is detected, the device sounds an alarm, giving people the chance to evacuate and ventilate the building. Unlike smoke detectors, carbon monoxide detectors do not need to be placed near ceiling level.
The use of carbon monoxide detectors has been standardized in many areas. In the US, NFPA 720–2009, the carbon monoxide detector guidelines published by the National Fire Protection Association, mandates the placement of carbon monoxide detectors/alarms on every level of the residence, including the basement, in addition to outside sleeping areas. In new homes, AC-powered detectors must have battery backup and be interconnected to ensure early warning of occupants at all levels. NFPA 720-2009 is the first national carbon monoxide standard to address devices in non-residential buildings. These guidelines, which now pertain to schools, healthcare centers, nursing homes, and other non-residential buildings, include three main points:
1. A secondary power supply (battery backup) must operate all carbon monoxide notification appliances for at least 12 hours,
2. Detectors must be on the ceiling in the same room as permanently installed fuel-burning appliances, and
3. Detectors must be located on every habitable level and in every HVAC zone of the building.
Gas organizations will often recommend getting gas appliances serviced at least once a year.
Legal requirements
The NFPA standard is not necessarily enforced by law. As of April 2006, the US state of Massachusetts requires detectors to be present in all residences with potential CO sources, regardless of building age and whether they are owner-occupied or rented. This is enforced by municipal inspectors and was inspired by the death of 7-year-old Nicole Garofalo in 2005 due to snow blocking a home heating vent. Other jurisdictions may have no requirement or only mandate detectors for new construction or at time of sale.
World Health Organization recommendations
The following guideline values (ppm values rounded) and periods of time-weighted average exposures have been determined in such a way that the carboxyhemoglobin (COHb) level of 2.5% is not exceeded, even when a normal subject engages in light or moderate exercise:
100 mg/m3 (87 ppm) for 15 min
60 mg/m3 (52 ppm) for 30 min
30 mg/m3 (26 ppm) for 1 h
10 mg/m3 (9 ppm) for 8 h
7 mg/m3 (6 ppm) for 24 h (for indoor air quality, so as not to exceed 2% COHb for chronic exposure)
Diagnosis
As many symptoms of carbon monoxide poisoning also occur with many other types of poisonings and infections (such as the flu), the diagnosis is often difficult. A history of potential carbon monoxide exposure, such as being exposed to a residential fire, may suggest poisoning, but the diagnosis is confirmed by measuring the levels of carbon monoxide in the blood. This can be determined by measuring the amount of carboxyhemoglobin compared to the amount of hemoglobin in the blood.
The ratio of carboxyhemoglobin to hemoglobin molecules in an average person may be up to 5%, although cigarette smokers who smoke two packs per day may have levels up to 9%. In symptomatic poisoned people they are often in the 10–30% range, while persons who die may have postmortem blood levels of 30–90%.
As people may continue to experience significant symptoms of CO poisoning long after their blood carboxyhemoglobin concentration has returned to normal, presenting to examination with a normal carboxyhemoglobin level (which may happen in late states of poisoning) does not rule out poisoning.
Measuring
Carbon monoxide may be quantitated in blood using spectrophotometric methods or chromatographic techniques in order to confirm a diagnosis of poisoning in a person or to assist in the forensic investigation of a case of fatal exposure.
A CO-oximeter can be used to determine carboxyhemoglobin levels. Pulse CO-oximeters estimate carboxyhemoglobin with a non-invasive finger clip similar to a pulse oximeter. These devices function by passing various wavelengths of light through the fingertip and measuring the light absorption of the different types of hemoglobin in the capillaries. The use of a regular pulse oximeter is not effective in the diagnosis of carbon monoxide poisoning as these devices may be unable to distinguish carboxyhemoglobin from oxyhemoglobin.
Breath CO monitoring offers an alternative to pulse CO-oximetry. Carboxyhemoglobin levels have been shown to have a strong correlation with breath CO concentration. However, many of these devices require the user to inhale deeply and hold their breath to allow the CO in the blood to escape into the lung before the measurement can be made. As this is not possible in people who are unresponsive, these devices may not appropriate for use in on-scene emergency care detection of CO poisoning.
Differential diagnosis
There are many conditions to be considered in the differential diagnosis of carbon monoxide poisoning. The earliest symptoms, especially from low level exposures, are often non-specific and readily confused with other illnesses, typically flu-like viral syndromes, depression, chronic fatigue syndrome, chest pain, and migraine or other headaches. Carbon monoxide has been called a "great mimicker" due to the presentation of poisoning being diverse and nonspecific. Other conditions included in the differential diagnosis include acute respiratory distress syndrome, altitude sickness, lactic acidosis, diabetic ketoacidosis, meningitis, methemoglobinemia, or opioid or toxic alcohol poisoning.
Treatment
Initial treatment for carbon monoxide poisoning is to immediately remove the person from the exposure without endangering further people. Those who are unconscious may require CPR on site. Administering oxygen via non-rebreather mask shortens the half-life of carbon monoxide from 320 minutes, when breathing normal air, to only 80 minutes. Oxygen hastens the dissociation of carbon monoxide from carboxyhemoglobin, thus turning it back into hemoglobin. Due to the possible severe effects in the baby, pregnant women are treated with oxygen for longer periods of time than non-pregnant people.
Hyperbaric oxygen
Hyperbaric oxygen is also used in the treatment of carbon monoxide poisoning, as it may hasten dissociation of CO from carboxyhemoglobin and cytochrome oxidase to a greater extent than normal oxygen. Hyperbaric oxygen at three times atmospheric pressure reduces the half life of carbon monoxide to 23 minutes, compared to 80 minutes for oxygen at regular atmospheric pressure. It may also enhance oxygen transport to the tissues by plasma, partially bypassing the normal transfer through hemoglobin. However, it is controversial whether hyperbaric oxygen actually offers any extra benefits over normal high flow oxygen, in terms of increased survival or improved long-term outcomes. There have been randomized controlled trials in which the two treatment options have been compared; of the six performed, four found hyperbaric oxygen improved outcome and two found no benefit for hyperbaric oxygen. Some of these trials have been criticized for apparent flaws in their implementation. A review of all the literature concluded that the role of hyperbaric oxygen is unclear and the available evidence neither confirms nor denies a medically meaningful benefit. The authors suggested a large, well designed, externally audited, multicentre trial to compare normal oxygen with hyperbaric oxygen. While hyperbaric oxygen therapy is used for severe poisonings, the benefit over standard oxygen delivery is unclear.
Other
Further treatment for other complications such as seizure, hypotension, cardiac abnormalities, pulmonary edema, and acidosis may be required. Hypotension requires treatment with intravenous fluids; vasopressors may be required to treat myocardial depression. Cardiac dysrhythmias are treated with standard advanced cardiac life support protocols. If severe, metabolic acidosis is treated with sodium bicarbonate. Treatment with sodium bicarbonate is controversial as acidosis may increase tissue oxygen availability. Treatment of acidosis may only need to consist of oxygen therapy. The delayed development of neuropsychiatric impairment is one of the most serious complications of carbon monoxide poisoning. Brain damage is confirmed following MRI or CAT scans. Extensive follow up and supportive treatment is often required for delayed neurological damage. Outcomes are often difficult to predict following poisoning, especially people who have symptoms of cardiac arrest, coma, metabolic acidosis, or have high carboxyhemoglobin levels. One study reported that approximately 30% of people with severe carbon monoxide poisoning will have a fatal outcome. It has been reported that electroconvulsive therapy (ECT) may increase the likelihood of delayed neuropsychiatric sequelae (DNS) after carbon monoxide (CO) poisoning. A device that also provides some carbon dioxide to stimulate faster breathing (sold under the brand name ClearMate) may also be used.
Pathophysiology
The precise mechanisms by which the effects of carbon monoxide are induced upon bodily systems are complex and not yet fully understood. Known mechanisms include carbon monoxide binding to hemoglobin, myoglobin and mitochondrial cytochrome c oxidase and restricting oxygen supply, and carbon monoxide causing brain lipid peroxidation.
Hemoglobin
Carbon monoxide has a higher diffusion coefficient compared to oxygen, and the main enzyme in the human body that produces carbon monoxide is heme oxygenase, which is located in nearly all cells and platelets. Most endogenously produced CO is stored bound to hemoglobin as carboxyhemoglobin. The simplistic understanding for the mechanism of carbon monoxide toxicity is based on excess carboxyhemoglobin decreasing the oxygen-delivery capacity of the blood to tissues throughout the body. In humans, the affinity between hemoglobin and carbon monoxide is approximately 240 times stronger than the affinity between hemoglobin and oxygen. However, certain mutations such as the Hb-Kirklareli mutation has a relative 80,000 times greater affinity for carbon monoxide than oxygen resulting in systemic carboxyhemoglobin reaching a sustained level of 16% COHb.
Hemoglobin is a tetramer with four prosthetic heme groups to serve as oxygen binding sites. The average red blood cell contains 250 million hemoglobin molecules, therefore 1 billion heme sites capable of binding gas. The binding of carbon monoxide at any one of these sites increases the oxygen affinity of the remaining three sites, which causes the hemoglobin molecule to retain oxygen that would otherwise be delivered to the tissue; therefore carbon monoxide binding at any site may be as dangerous as carbon monoxide binding to all sites. Delivery of oxygen is largely driven by the Bohr effect and Haldane effect. To provide a simplified synopsis of the molecular mechanism of systemic gas exchange in layman's terms, upon inhalation of air it was widely thought oxygen binding to any of the heme sites triggers a conformational change in the globin/protein unit of hemoglobin which then enables the binding of additional oxygen to each of the other vacant heme sites. Upon arrival to the cell/tissues, oxygen release into the tissue is driven by "acidification" of the local pH (meaning a relatively higher concentration of 'acidic' protons/hydrogen ions) caused by an increase in the biotransformation of carbon dioxide waste into carbonic acid via carbonic anhydrase. In other words, oxygenated arterial blood arrives at cells in the "hemoglobin R-state" which has deprotonated/unionized amino acid residues (regarding nitrogen/amines) due to the less-acidic arterial pH environment (arterial blood averages pH 7.407 whereas venous blood is slightly more acidic at pH 7.371). The "T-state" of hemoglobin is deoxygenated in venous blood partially due to protonation/ionization caused by the acidic environment hence causing a conformation unsuited for oxygen-binding (in other words, oxygen is 'ejected' upon arrival to the cell because acid "attacks" the amines of hemoglobin causing ionization/protonation of the amine residues resulting in a conformation change unsuited for retaining oxygen). Furthermore, the mechanism for formation of carbaminohemoglobin generates additional 'acidic' hydrogen ions that may further stabilize the protonated/ionized deoxygenated hemoglobin. Upon return of venous blood into the lung and subsequent exhalation of carbon dioxide, the blood is "de-acidified" (see also: hyperventilation) allowing for the deprotonation/unionization of hemoglobin to then re-enable oxygen-binding as part of the transition to arterial blood (note this process is complex due to involvement of chemoreceptors and other physiological functionalities). Carbon monoxide is not 'ejected' due to acid, therefore carbon monoxide poisoning disturbs this physiological process hence the venous blood of poisoning patients is bright red akin to arterial blood since the carbonyl/carbon monoxide is retained. Hemoglobin is dark in deoxygenated venous blood, but it has a bright red color when carrying blood in oxygenated arterial blood and when converted into carboxyhemoglobin in both arterial and venous blood, so poisoned cadavers and even commercial meats treated with carbon monoxide acquire an unnatural lively reddish hue.
At toxic concentrations, carbon monoxide as carboxyhemoglobin significantly interferes with respiration and gas exchange by simultaneously inhibiting acquisition and delivery of oxygen to cells and preventing formation of carbaminohemoglobin which accounts for approximately 30% of carbon dioxide exportation. Therefore, a patient with carbon monoxide poisoning may experience severe hypoxia and acidosis (potentially both respiratory acidosis and metabolic acidosis) in addition to the toxicities of excess carbon monoxide inhibiting numerous hemoproteins, metallic and non-metallic targets which affect cellular machinery.
Myoglobin
Carbon monoxide also binds to the hemeprotein myoglobin. It has a high affinity for myoglobin, about 60 times greater than that of oxygen. Carbon monoxide bound to myoglobin may impair its ability to utilize oxygen. This causes reduced cardiac output and hypotension, which may result in brain ischemia. A delayed return of symptoms have been reported. This results following a recurrence of increased carboxyhemoglobin levels; this effect may be due to a late release of carbon monoxide from myoglobin, which subsequently binds to hemoglobin.
Cytochrome oxidase
Another mechanism involves effects on the mitochondrial respiratory enzyme chain that is responsible for effective tissue utilization of oxygen. Carbon monoxide binds to cytochrome oxidase with less affinity than oxygen, so it is possible that it requires significant intracellular hypoxia before binding. This binding interferes with aerobic metabolism and efficient adenosine triphosphate synthesis. Cells respond by switching to anaerobic metabolism, causing anoxia, lactic acidosis, and eventual cell death. The rate of dissociation between carbon monoxide and cytochrome oxidase is slow, causing a relatively prolonged impairment of oxidative metabolism.
Central nervous system effects
The mechanism that is thought to have a significant influence on delayed effects involves formed blood cells and chemical mediators, which cause brain lipid peroxidation (degradation of unsaturated fatty acids). Carbon monoxide causes endothelial cell and platelet release of nitric oxide, and the formation of oxygen free radicals including peroxynitrite. In the brain this causes further mitochondrial dysfunction, capillary leakage, leukocyte sequestration, and apoptosis. The result of these effects is lipid peroxidation, which causes delayed reversible demyelination of white matter in the central nervous system known as Grinker myelinopathy, which can lead to edema and necrosis within the brain. This brain damage occurs mainly during the recovery period. This may result in cognitive defects, especially affecting memory and learning, and movement disorders. These disorders are typically related to damage to the cerebral white matter and basal ganglia. Hallmark pathological changes following poisoning are bilateral necrosis of the white matter, globus pallidus, cerebellum, hippocampus and the cerebral cortex.
Pregnancy
Carbon monoxide poisoning in pregnant women may cause severe adverse fetal effects. Poisoning causes fetal tissue hypoxia by decreasing the release of maternal oxygen to the fetus. Carbon monoxide also crosses the placenta and combines with fetal hemoglobin, causing more direct fetal tissue hypoxia. Additionally, fetal hemoglobin has a 10 to 15% higher affinity for carbon monoxide than adult hemoglobin, causing more severe poisoning in the fetus than in the adult. Elimination of carbon monoxide is slower in the fetus, leading to an accumulation of the toxic chemical. The level of fetal morbidity and mortality in acute carbon monoxide poisoning is significant, so despite mild maternal poisoning or following maternal recovery, severe fetal poisoning or death may still occur.
History
Humans have maintained a complex relationship with carbon monoxide since first learning to control fire circa 800,000 BC. Primitive cavemen probably discovered the toxicity of carbon monoxide upon introducing fire into their dwellings. The early development of metallurgy and smelting technologies emerging circa 6,000 BC through the Bronze Age likewise plagued humankind with carbon monoxide exposure. Apart from the toxicity of carbon monoxide, indigenous Native Americans may have experienced the neuroactive properties of carbon monoxide through shamanistic fireside rituals.
Early civilizations developed mythological tales to explain the origin of fire, such as Vulcan, Pkharmat, and Prometheus from Greek mythology who shared fire with humans. Aristotle (384–322 BC) first recorded that burning coals produced toxic fumes. Greek physician Galen (129–199 AD) speculated that there was a change in the composition of the air that caused harm when inhaled, and symptoms of CO poisoning appeared in Cassius Iatrosophista's Quaestiones Medicae et Problemata Naturalia circa 130 AD. Julian the Apostate, Caelius Aurelianus, and several others similarly documented early knowledge of the toxicity symptoms of carbon monoxide poisoning as caused by coal fumes in the ancient era.
Documented cases by Livy and Cicero allude to carbon monoxide being used as a method of suicide in ancient Rome. Emperor Lucius Verus used smoke to execute prisoners. Many deaths have been linked to carbon monoxide poisoning including Emperor Jovian, Empress Fausta, and Seneca. The most high-profile death by carbon monoxide poisoning may possibly have been Cleopatra or Edgar Allan Poe.
In the fifteenth century, coal miners believed sudden death was caused by evil spirits; carbon monoxide poisoning has been linked to supernatural and paranormal experiences, witchcraft, etc. throughout the following centuries including in the modern present day exemplified by Carrie Poppy's investigations.
Georg Ernst Stahl mentioned carbonarii halitus in 1697 in reference to toxic vapors thought to be carbon monoxide. Friedrich Hoffmann conducted the first modern scientific investigation into carbon monoxide poisoning from coal in 1716, notably rejecting villagers attributing death to demonic superstition. Herman Boerhaave conducted the first scientific experiments on the effect of carbon monoxide (coal fumes) on animals in the 1730s. Joseph Priestley is credited with first synthesizing carbon monoxide in 1772 which he had called heavy inflammable air, and Carl Wilhelm Scheele isolated carbon monoxide from coal in 1773 suggesting it to be the toxic entity.
The dose-dependent risk of carbon monoxide poisoning as hydrocarbonate was investigated in the late 1790s by Thomas Beddoes, James Watt, Tiberius Cavallo, James Lind, Humphry Davy, and many others in the context of inhalation of factitious airs, much of which occurred at the Pneumatic Institution.
William Cruickshank discovered carbon monoxide as a molecule containing one carbon and one oxygen atom in 1800, thereby initiating the modern era of research exclusively focused on carbon monoxide. The mechanism for toxicity was first suggested by James Watt in 1793, followed by Adrien Chenot in 1854 and finally demonstrated by Claude Bernard after 1846 as published in 1857 and also independently published by Felix Hoppe-Seyler in the same year.
The first controlled clinical trial studying the toxicity of carbon monoxide occurred in 1973.
Historical detection
Carbon monoxide poisoning has plagued coal miners for many centuries. In the context of mining, carbon monoxide is widely known as whitedamp. John Scott Haldane identified carbon monoxide as the lethal constituent of afterdamp, the gas created by combustion, after examining many bodies of miners killed in pit explosions. By 1911, Haldane introduced the use of small animals for miners to detect dangerous levels of carbon monoxide underground, either white mice or canaries which have little tolerance for carbon monoxide thereby offering an early warning, i.e. canary in a coal mine. The canary in British pits was replaced in 1986 by the electronic gas detector.
The first qualitative analytical method to detect carboxyhemoglobin emerged in 1858 with a colorimetric method developed by Felix Hoppe-Seyler, and the first quantitative analysis method emerged in 1880 with Josef von Fodor.
Historical treatment
The use of oxygen emerged with anecdotal reports such as Humphry Davy having been treated with oxygen in 1799 upon inhaling three quarts of hydrocarbonate (water gas). Samuel Witter developed an oxygen inhalation protocol in response to carbon monoxide poisoning in 1814. Similarly, an oxygen inhalation protocol was recommend for malaria (literally translated to "bad air") in 1830 based on malaria symptoms aligning with carbon monoxide poisoning. Other oxygen protocols emerged in the late 1800s. The use of hyperbaric oxygen in rats following poisoning was studied by Haldane in 1895 while its use in humans began in the 1960s.
Incidents
The worst accidental mass poisoning from carbon monoxide was the Balvano train disaster which occurred on 3 March 1944 in Italy, when a freight train with many illegal passengers stalled in a tunnel, leading to the death of over 500 people.
Over 50 people are suspected to have died from smoke inhalation as a result of the Branch Davidian Massacre during the Waco siege in 1993.
On 14 December 2024 12 individuals died by carbon monoxide poisoning in Gudauri (Georgia) as electric generators using fuel oil were placed in a closed area near their rooms.
Weaponization
In ancient history, Hannibal executed Roman prisoners with coal fumes during the Second Punic War.
The extermination of stray dogs by a carbon monoxide gas chamber was described in 1874. In 1884, an article appeared in Scientific American describing the use of a carbon monoxide gas chamber for slaughterhouse operations as well as euthanizing a variety of animals.
As part of the Holocaust during World War II, the Nazis used gas vans at Chelmno extermination camp and elsewhere to murder an estimated 700,000 or more people by carbon monoxide poisoning. This method was also used in the gas chambers of several death camps such as Treblinka, Sobibor, and Belzec. Gassing with carbon monoxide started in Action T4. The gas was supplied by IG Farben in pressurized cylinders and fed by tubes into the gas chambers built at various mental hospitals, such as Hartheim Euthanasia Centre. Exhaust fumes from tank engines, for example, were used to supply the gas to the chambers.
| Biology and health sciences | Types | Health |
488748 | https://en.wikipedia.org/wiki/Reionization | Reionization | In the fields of Big Bang theory and cosmology, reionization is the process that caused electrically neutral atoms in the universe to reionize after the lapse of the "dark ages".
Detecting and studying the reionization process is challenging but multiple avenues have been pursued.
This reionization was driven by the formation of the first stars and galaxies.
Concept
Reionization refers to a change in the intergalactic medium from neutral hydrogen to ions. The neutral hydrogen had been ions at an earlier stage in the history of the universe, thus the conversion back into ions is termed a reionization. The reionization was driven by energetic photons emitted by the first stars and galaxies.
In the timeline of the universe, neutral hydrogen gas was originally formed when primordial hydrogen nuclei (protons) combined with electrons. Light with sufficient energy will ionize neutral hydrogen gas. At early times, light was so dense and energetic that hydrogen gas was immediately ionized. As the universe expanded and cooled, the rate of recombination of electrons and protons to form neutral hydrogen was higher than the ionization rate. At around 379,000 years after the Big Bang (redshift z = 1089), this recombination left most normal matter in the form of neutral hydrogen.
The universe was opaque before the recombination, due to the scattering of photons of all wavelengths off free electrons (and free protons, to a significantly lesser extent), but it became increasingly transparent as more electrons and protons combined to form neutral hydrogen atoms. While the electrons of neutral hydrogen can absorb photons of some wavelengths by rising to an excited state, a universe full of neutral hydrogen will be relatively opaque only at those few wavelengths. The remaining light could travel freely and become the cosmic microwave background radiation. The only other light at this point would be provided by those excited hydrogen atoms, marking the beginning of an era called the Dark Ages of the universe.
The second phase change occurred once gas clouds started to condense in the early universe that were energetic enough to re-ionize neutral hydrogen. As these objects formed and radiated energy, the universe reverted from being composed of neutral atoms, to once again being an ionized plasma. This occurred between 150 million and one billion years after the Big Bang (at a redshift 20 > z > 6) At that time, however, matter had been diffused by the expansion of the universe, and the scattering interactions of photons and electrons were much less frequent than before electron-proton recombination. Thus, the universe was full of low density ionized hydrogen and remained transparent, as is the case today.
It is believed that the primordial helium also experienced a similar reionization phase change, but at a later epoch in the history of the universe.
Stages
Theoretical models give a timeline of the reionization process.
In the first stage of reionization, each new star is surrounded by neutral hydrogen. Light emitted by the star ionizes gas immediately around the star. Then light can reach further out to ionize gas. The ions can recombine, competing with the ionization process. The ionized gas will be hot and it will expand, clearing out the region around the star. The sphere of ionized gas expands until the amount of light from the star that can cause ionizations balances the recombination, a process that takes hundreds of millions of years. (The time is so long that stars die before the full extent of the reionization completes for that star.) At some point the shell of ionization from each star in a galaxy begin to overlap and the ionization frontier pushes out into the intergalatic medium.
Detection methods
Looking back so far in the history of the universe presents some observational challenges. There are, however, a few observational methods for studying reionization.
Quasars and the Gunn-Peterson trough
One means of studying reionization uses the spectra of distant quasars. Quasars release an extraordinary amount of energy, being among the brightest objects in the universe. As a result, some quasars are detectable from as long ago as the epoch of reionization. Quasars also happen to have relatively uniform spectral features, regardless of their position in the sky or distance from the Earth. Thus it can be inferred that any major differences between quasar spectra will be caused by the interaction of their emission with atoms along the line of sight. For wavelengths of light at the energies of one of the Lyman transitions of hydrogen, the scattering cross-section is large, meaning that even for low levels of neutral hydrogen in the intergalactic medium (IGM), absorption at those wavelengths is highly likely.
For nearby objects in the universe, spectral absorption lines are very sharp, as only photons with energies just sufficient to cause an atomic transition can cause that transition. However, the distances between quasars and the telescopes which detect them are large, which means that the expansion of the universe causes light to undergo noticeable redshifting. This means that as light from the quasar travels through the IGM and is redshifted, wavelengths which had been below the Lyman Alpha limit are stretched, and will in effect begin to fill in the Lyman absorption band. This means that instead of showing sharp spectral absorption lines, a quasar's light which has traveled through a large, spread out region of neutral hydrogen will show a Gunn-Peterson trough.
The redshifting for a particular quasar provides temporal information about reionization. Since an object's redshift corresponds to the time at which it emitted the light, it is possible to determine when reionization ended. Quasars below a certain redshift (closer in space and time) do not show the Gunn-Peterson trough (though they may show the Lyman-alpha forest), while quasars emitting light prior to reionization will feature a Gunn-Peterson trough. In 2001, four quasars were detected by the Sloan Digital Sky Survey with redshifts ranging from z = 5.82 to z = 6.28. While the quasars above z = 6 showed a Gunn-Peterson trough, indicating that the IGM was still at least partly neutral, the ones below did not, meaning the hydrogen was ionized. As reionization is expected to occur over relatively short timescales, the results suggest that the universe was approaching the end of reionization at z = 6. This, in turn, suggests that the universe must still have been almost entirely neutral at z > 10. On the other hand, long absorption troughs persisting down to z < 5.5 in the Lyman-alpha and Lyman-beta forests suggest that reionization potentially extends later than z = 6.
CMB anisotropy and polarization
The anisotropy of the cosmic microwave background on different angular scales can also be used to study reionization. Photons undergo scattering when there are free electrons present, in a process known as Thomson scattering. However, as the universe expands, the density of free electrons will decrease, and scattering will occur less frequently. In the period during and after reionization, but before significant expansion had occurred to sufficiently lower the electron density, the light that composes the CMB will experience observable Thomson scattering. This scattering will leave its mark on the CMB anisotropy map, introducing secondary anisotropies (anisotropies introduced after recombination). The overall effect is to erase anisotropies that occur on smaller scales. While anisotropies on small scales are erased, polarization anisotropies are actually introduced because of reionization. By looking at the CMB anisotropies observed, and comparing with what they would look like had reionization not taken place, the electron column density at the time of reionization can be determined. With this, the age of the universe when reionization occurred can then be calculated.
The Wilkinson Microwave Anisotropy Probe allowed that comparison to be made. The initial observations, released in 2003, suggested that reionization took place from 30 > z > 11. This redshift range was in clear disagreement with the results from studying quasar spectra. However, the three year WMAP data returned a different result, with reionization beginning at z = 11 and the universe ionized by z = 7. This is in much better agreement with the quasar data.
Results in 2018 from Planck mission, yield an instantaneous reionization redshift of z = 7.68 ± 0.79.
The parameter usually quoted here is τ, the "optical depth to reionization," or alternatively, zre, the redshift of reionization, assuming it was an instantaneous event. While this is unlikely to be physical, since reionization was very likely not instantaneous, zre provides an estimate of the mean redshift of reionization.
Lyman alpha emission
Lyman alpha light from galaxies offers a complementary tool set to study reionization. The Lyman alpha line is the n=2 to n=1 transition of neutral hydrogen, and can be produced copiously by galaxies with young stars. Moreover, Lyman alpha photons interact strongly with neutral hydrogen in intergalactic gas through resonant scattering, wherein neutral atoms in the ground (n=1) state absorb Lyman alpha photons and almost immediately re-emit them in a random direction. This obscures Lyman alpha emission from galaxies that are embedded in neutral gas. Thus, experiments to find galaxies by their Lyman alpha light can indicate the ionization state of the surrounding gas. An average density of galaxies with detectable Lyman alpha emission means the surrounding gas must be ionized; while an absence of detectable Lyman alpha sources may indicate neutral regions. A closely related class of experiments measures the Lyman alpha line strength in samples of galaxies identified by other methods (primarily Lyman break galaxy searches).
The earliest application of this method was in 2004, when the tension between late neutral gas indicated by quasar spectra and early reionization suggested by CMB results was strong. The detection of Lyman alpha galaxies at redshift z=6.5 demonstrated that the intergalactic gas was already predominantly ionized at an earlier time than the quasar spectra suggested. Subsequent applications of the method suggested some residual neutral gas as recently as z=6.5, but still indicate that a majority of intergalactic gas was ionized prior to z=7.
Lyman alpha emission can be used in other ways to further probe reionization. Theory suggests that reionization was patchy, meaning that the clustering of Lyman alpha selected samples should be strongly enhanced during the middle phases of reionization. Moreover, specific ionized regions can be pinpointed by identifying groups of Lyman alpha emitters.
21-cm line
Even with the quasar data roughly in agreement with the CMB anisotropy data, there are still a number of questions, especially concerning the energy sources of reionization and the effects on, and role of, structure formation during reionization. The 21-cm line in hydrogen is potentially a means of studying this period, as well as the "dark ages" that preceded reionization. The 21-cm line occurs in neutral hydrogen, due to differences in energy between the spin triplet and spin singlet states of the electron and proton. This transition is forbidden, meaning it occurs extremely rarely. The transition is also highly temperature dependent, meaning that as objects form in the "dark ages" and emit Lyman-alpha photons that are absorbed and re-emitted by surrounding neutral hydrogen, it will produce a 21-cm line signal in that hydrogen through Wouthuysen-Field coupling. By studying 21-cm line emission, it will be possible to learn more about the early structures that formed. Observations from the Experiment to Detect the Global Epoch of Reionization Signature (EDGES) points to a signal from this era, although follow-up observations will be needed to confirm it. Several other projects hope to make headway in this area in the near future, such as the Precision Array for Probing the Epoch of Reionization (PAPER), Low Frequency Array (LOFAR), Murchison Widefield Array (MWA), Giant Metrewave Radio Telescope (GMRT), Mapper of the IGM Spin Temperature (MIST), the Dark Ages Radio Explorer (DARE) mission, and the Large-Aperture Experiment to Detect the Dark Ages (LEDA).
Energy sources
While observations have come in which narrow the window during which the epoch of reionization could have taken place, it is still uncertain which objects provided the photons that reionized the IGM. To ionize neutral hydrogen, an energy larger than 13.6 eV is required, which corresponds to photons with a wavelength of 91.2 nm or shorter. This is in the ultraviolet part of the electromagnetic spectrum, which means that the primary candidates are all sources which produce a significant amount of energy in the ultraviolet and above. How numerous the source is must also be considered, as well as the longevity, as protons and electrons will recombine if energy is not continuously provided to keep them apart. Altogether, the critical parameter for any source considered can be summarized as its "emission rate of hydrogen-ionizing photons per unit cosmological volume." With these constraints, it is expected that quasars and first generation stars and galaxies were the main sources of energy.
Dwarf galaxies
Dwarf galaxies are currently considered to be the primary source of ionizing photons during the epoch of reionization. For most scenarios, this would require the log-slope of the UV galaxy luminosity function, often denoted α, to be steeper than it is today, approaching α = -2. With the advent of the James Webb Space Telescope (JWST), constraints on the UV luminosity function at the Epoch of Reionization have become commonplace, allowing for better constraints on the faint, low-mass population of galaxies.
In 2014, two separate studies identified two Green Pea galaxies (GPs) to be likely Lyman Continuum (LyC)-emitting candidates. Compact dwarf star-forming galaxies like the GPs are considered excellent low-redshift analogs of high-redshift Lyman-alpha and LyC emitters (LAEs and LCEs, respectively). At that time, only two other LCEs were known: Haro 11 and Tololo-1247-232. Finding local LyC emitters has thus become crucial to the theories about the early universe and the epoch of reionization.
Subsequently, motivated, a series of surveys have been conducted using Hubble Space Telescope's Cosmic Origins Spectrograph (HST/COS) to measure the LyC directly. These efforts culminated in the Low-redshift Lyman Continuum Survey, a large HST/COS program which nearly tripled the number of direct measurements of the LyC from dwarf galaxies. To date, at least 50 LCEs have been confirmed using HST/COS with LyC escape fractions anywhere from ≈ 0 to 88%. The results from the Low-redshift Lyman Continuum Survey have provided the empirical foundation necessary to identify and understand LCEs at the Epoch of Reionization. With new observations from JWST, populations of LCEs are now being studied at cosmological redshifts greater than 6, allowing for the first time a detailed and direct assessment of the origins of cosmic Reionization. Combining these large samples of galaxies with new constraints on the UV luminosity function indicates that dwarf galaxies overwhelmingly contribute to Reionization.
Quasars
Quasars, a class of active galactic nuclei (AGN), were considered a good candidate source because they are highly efficient at converting mass to energy, and emit a great deal of light above the threshold for ionizing hydrogen. It is unknown, however, how many quasars existed prior to reionization. Only the brightest of quasars present during reionization can be detected, which means there is no direct information about dimmer quasars that existed. However, by looking at the more easily observed quasars in the nearby universe, and assuming that the luminosity function (number of quasars as a function of luminosity) during reionization will be approximately the same as it is today, it is possible to make estimates of the quasar populations at earlier times. Such studies have found that quasars do not exist in high enough numbers to reionize the IGM alone, saying that "only if the ionizing background is dominated by low-luminosity AGNs can the quasar luminosity function provide enough ionizing photons."
Population III stars
Population III stars were the earliest stars, which had no elements more massive than hydrogen or helium. During Big Bang nucleosynthesis, the only elements that formed aside from hydrogen and helium were trace amounts of lithium. Yet quasar spectra have revealed the presence of heavy elements in the intergalactic medium at an early era. Supernova explosions produce such heavy elements, so hot, large, Population III stars which will form supernovae are a possible mechanism for reionization. While they have not been directly observed, they are consistent according to models using numerical simulation and current observations. A gravitationally lensed galaxy also provides indirect evidence of Population III stars. Even without direct observations of Population III stars, they are a compelling source. They are more efficient and effective ionizers than Population II stars, as they emit more ionizing photons, and are capable of reionizing hydrogen on their own in some reionization models with reasonable initial mass functions. As a consequence, Population III stars are currently considered the most likely energy source to initiate the reionization of the universe, though other sources are likely to have taken over and driven reionization to completion.
In June 2015, astronomers reported evidence for Population III stars in the Cosmos Redshift 7 galaxy at . Such stars are likely to have existed in the very early universe (i.e., at high redshift), and may have started the production of chemical elements heavier than hydrogen that are needed for the later formation of planets and life as we know it.
| Physical sciences | Physical cosmology | Astronomy |
488815 | https://en.wikipedia.org/wiki/Utility%20frequency | Utility frequency | The utility frequency, (power) line frequency (American English) or mains frequency (British English) is the nominal frequency of the oscillations of alternating current (AC) in a wide area synchronous grid transmitted from a power station to the end-user. In large parts of the world this is 50 Hz, although in the Americas and parts of Asia it is typically 60 Hz. Current usage by country or region is given in the list of mains electricity by country.
During the development of commercial electric power systems in the late-19th and early-20th centuries, many different frequencies (and voltages) had been used. Large investment in equipment at one frequency made standardization a slow process. However, as of the turn of the 21st century, places that now use the 50 Hz frequency tend to use 220–240 V, and those that now use 60 Hz tend to use 100–127 V. Both frequencies coexist today (Japan uses both) with no great technical reason to prefer one over the other and no apparent desire for complete worldwide standardization.
Electric clocks
In practice, the exact frequency of the grid varies around the nominal frequency, reducing when the grid is heavily loaded, and speeding up when lightly loaded. However, most utilities will adjust generation onto the grid over the course of the day to ensure a constant number of cycles occur. This is used by some clocks to accurately maintain their time.
Operating factors
Several factors influence the choice of frequency in an AC system. Lighting, motors, transformers, generators, and transmission lines all have characteristics which depend on the power frequency. All of these factors interact and make selection of a power frequency a matter of considerable importance. The best frequency is a compromise among competing requirements.
In the late 19th century, designers would pick a relatively high frequency for systems featuring transformers and arc lights, so as to economize on transformer materials and to reduce visible flickering of the lamps, but would pick a lower frequency for systems with long transmission lines or feeding primarily motor loads or rotary converters for producing direct current. When large central generating stations became practical, the choice of frequency was made based on the nature of the intended load. Eventually improvements in machine design allowed a single frequency to be used both for lighting and motor loads. A unified system improved the economics of electricity production, since system load was more uniform during the course of a day.
Lighting
The first applications of commercial electric power were incandescent lighting and commutator-type electric motors. Both devices operate well on DC, but DC could not be easily changed in voltage, and was generally only produced at the required utilization voltage.
If an incandescent lamp is operated on a low-frequency current, the filament cools on each half-cycle of the alternating current, leading to perceptible change in brightness and flicker of the lamps; the effect is more pronounced with arc lamps, and the later mercury-vapor lamps and fluorescent lamps. Open arc lamps made an audible buzz on alternating current, leading to experiments with high-frequency alternators to raise the sound above the range of human hearing.
Rotating machines
Commutator-type motors do not operate well on high-frequency AC, because the rapid changes of current are opposed by the inductance of the motor field. Though commutator-type universal motors are common in AC household appliances and power tools, they are small motors, less than 1 kW. The induction motor was found to work well on frequencies around 50 to 60 Hz, but with the materials available in the 1890s would not work well at a frequency of, say, 133 Hz. There is a fixed relationship between the number of magnetic poles in the induction motor field, the frequency of the alternating current, and the rotation speed; so, a given standard speed limits the choice of frequency (and the reverse). Once AC electric motors became common, it was important to standardize frequency for compatibility with the customer's equipment.
Generators operated by slow-speed reciprocating engines will produce lower frequencies, for a given number of poles, than those operated by, for example, a high-speed steam turbine. For very slow prime mover speeds, it would be costly to build a generator with enough poles to provide a high AC frequency. As well, synchronizing two generators to the same speed was found to be easier at lower speeds. While belt drives were common as a way to increase speed of slow engines, in very large ratings (thousands of kilowatts) these were expensive, inefficient, and unreliable. After about 1906, generators driven directly by steam turbines favored higher frequencies. The steadier rotation speed of high-speed machines allowed for satisfactory operation of commutators in rotary converters.
The synchronous speed N in RPM is calculated using the formula,
where f is the frequency in hertz and P is the number of poles.
Direct-current power was not entirely displaced by alternating current and was useful in railway and electrochemical processes. Prior to the development of mercury arc valve rectifiers, rotary converters were used to produce DC power from AC. Like other commutator-type machines, these worked better with lower frequencies.
Transmission and transformers
With AC, transformers can be used to step down high transmission voltages to lower customer utilization voltage. The transformer is effectively a voltage conversion device with no moving parts and requiring little maintenance. The use of AC eliminated the need for spinning DC voltage conversion motor-generators that require regular maintenance and monitoring.
Since, for a given power level, the dimensions of a transformer are roughly inversely proportional to frequency, a system with many transformers would be more economical at a higher frequency.
Electric power transmission over long lines favors lower frequencies. The effects of the distributed capacitance and inductance of the line are less at low frequency.
System interconnection
Generators can only be interconnected to operate in parallel if they are of the same frequency and wave-shape. By standardizing the frequency used, generators in a geographic area can be interconnected in a grid, providing reliability and cost savings.
History
Many different power frequencies were used in the 19th century.
Very early isolated AC generating schemes used arbitrary frequencies based on convenience for steam engine, water turbine, and electrical generator design. Frequencies between Hz and Hz were used on different systems. For example, the city of Coventry, England, in 1895 had a unique 87 Hz single-phase distribution system that was in use until 1906. The proliferation of frequencies grew out of the rapid development of electrical machines in the period 1880 through 1900.
In the early incandescent lighting period, single-phase AC was common and typical generators were 8-pole machines operated at 2,000 RPM, giving a frequency of 133 hertz.
Though many theories exist, and quite a few entertaining urban legends, there is little certitude in the details of the history of 60 Hz vs. 50 Hz.
The German company AEG (descended from a company founded by Edison in Germany) built the first German generating facility to run at 50 Hz. At the time, AEG had a virtual monopoly and their standard spread to the rest of Europe. After observing flicker of lamps operated by the 40 Hz power transmitted by the Lauffen-Frankfurt link in 1891, AEG raised their standard frequency to 50 Hz in 1891.
Westinghouse Electric decided to standardize on a higher frequency to permit operation of both electric lighting and induction motors on the same generating system. Although 50 Hz was suitable for both, in 1890 Westinghouse considered that existing arc-lighting equipment operated slightly better on 60 Hz, and so that frequency was chosen. The operation of Tesla's induction motor, licensed by Westinghouse in 1888, required a lower frequency than the 133 Hz common for lighting systems at that time. In 1893 General Electric Corporation, which was affiliated with AEG in Germany, built a generating project at Mill Creek to bring electricity to Redlands, California using 50 Hz, but changed to 60 Hz a year later to maintain market share with the Westinghouse standard.
25 Hz origins
The first generators at the Niagara Falls project, built by Westinghouse in 1895, were 25 Hz, because the turbine speed had already been set before alternating current power transmission had been definitively selected. Westinghouse would have selected a low frequency of 30 Hz to drive motor loads, but the turbines for the project had already been specified at 250 RPM. The machines could have been made to deliver Hz power suitable for heavy commutator-type motors, but the Westinghouse company objected that this would be undesirable for lighting and suggested Hz. Eventually a compromise of 25 Hz, with 12-pole 250 RPM generators, was chosen. Because the Niagara project was so influential on electric power systems design, 25 Hz prevailed as the North American standard for low-frequency AC.
40 Hz origins
A General Electric study concluded that 40 Hz would have been a good compromise between lighting, motor, and transmission needs, given the materials and equipment available in the first quarter of the 20th century. Several 40 Hz systems were built. The Lauffen-Frankfurt demonstration used 40 Hz to transmit power 175 km in 1891. A large interconnected 40 Hz network existed in north-east England (the Newcastle-upon-Tyne Electric Supply Company, NESCO) until the advent of the National Grid (UK) in the late 1920s, and projects in Italy used 42 Hz. The oldest continuously operating commercial hydroelectric power station in the United States, Mechanicville Hydroelectric Plant, still produces electric power at 40 Hz and supplies power to the local 60 Hz transmission system through frequency changers. Industrial plants and mines in North America and Australia sometimes were built with 40 Hz electrical systems which were maintained until too uneconomic to continue. Although frequencies near 40 Hz found much commercial use, these were bypassed by standardized frequencies of 25, 50 and 60 Hz preferred by higher volume equipment manufacturers.
The Ganz Company of Hungary had standardized on 5000 alternations per minute (41 Hz) for their products, so Ganz clients had 41 Hz systems that in some cases ran for many years.
Standardization
In the early days of electrification, so many frequencies were used that no single value prevailed (London in 1918 had ten different frequencies). As the 20th century continued, more power was produced at 60 Hz (North America) or 50 Hz (Europe and most of Asia). Standardization allowed international trade in electrical equipment. Much later, the use of standard frequencies allowed interconnection of power grids. It was not until after World War II – with the advent of affordable electrical consumer goods – that more uniform standards were enacted.
In the United Kingdom, a standard frequency of 50 Hz was declared as early as 1904, but significant development continued at other frequencies. The implementation of the National Grid starting in 1926 compelled the standardization of frequencies among the many interconnected electrical service providers. The 50 Hz standard was completely established only after World War II.
By about 1900, European manufacturers had mostly standardized on 50 Hz for new installations. The German Verband der Elektrotechnik (VDE), in the first standard for electrical machines and transformers in 1902, recommended 25 Hz and 50 Hz as standard frequencies. VDE did not see much application of 25 Hz, and dropped it from the 1914 edition of the standard. Remnant installations at other frequencies persisted until well after the Second World War.
Because of the cost of conversion, some parts of the distribution system may continue to operate on original frequencies even after a new frequency is chosen. 25 Hz power was used in Ontario, Quebec, the northern United States, and for railway electrification. In the 1950s, many 25 Hz systems, from the generators right through to household appliances, were converted and standardized. Until 2006, some 25 Hz generators were still in existence at the Sir Adam Beck 1 (these were retrofitted to 60 Hz) and the Rankine generating stations (until its 2006 closure) near Niagara Falls to provide power for large industrial customers who did not want to replace existing equipment; and some 25 Hz motors and a 25 Hz power station exist in New Orleans for floodwater pumps. The 15 kV AC rail networks, used in Germany, Austria, Switzerland, Sweden, and Norway, still operate at Hz or 16.7 Hz.
In some cases, where most load was to be railway or motor loads, it was considered economic to generate power at 25 Hz and install rotary converters for 60 Hz distribution. Converters for production of DC from alternating current were available in larger sizes and were more efficient at 25 Hz compared with 60 Hz. Remnant fragments of older systems may be tied to the standard frequency system via a rotary converter or static inverter frequency changer. These allow energy to be interchanged between two power networks at different frequencies, but the systems are large, costly, and waste some energy in operation.
Rotating-machine frequency changers used to convert between 25 Hz and 60 Hz systems were awkward to design; a 60 Hz machine with 24 poles would turn at the same speed as a 25 Hz machine with 10 poles, making the machines large, slow-speed, and expensive. A ratio of 60/30 would have simplified these designs, but the installed base at 25 Hz was too large to be economically opposed.
In the United States, Southern California Edison had standardized on 50 Hz. Much of Southern California operated on 50 Hz and did not completely change frequency of their generators and customer equipment to 60 Hz until around 1948. Some projects by the Au Sable Electric Company used 30 Hz at transmission voltages up to 110,000 volts in 1914.
Initially in Brazil, electric machinery were imported from Europe and United States, implying the country had both 50 Hz and 60 Hz standards according to each region. In 1938, the federal government made a law, Decreto-Lei 852, intended to bring the whole country under 50 Hz within eight years. The law did not work, and in the early 1960s it was decided that Brazil would be unified under 60 Hz standard, because most developed and industrialized areas used 60 Hz; and a new law Lei 4.454 was declared in 1964. Brazil underwent a frequency conversion program to 60 Hz that was not completed until 1978.
In Mexico, areas operating on 50 Hz grid were converted during the 1970s, uniting the country under 60 Hz.
In Japan, the western part of the country (Nagoya and west) uses 60 Hz and the eastern part (Tokyo and east) uses 50 Hz. This originates in the first purchases of generators from AEG in 1895, installed for Tokyo, and General Electric in 1896, installed in Osaka. The boundary between the two regions contains four back-to-back HVDC substations which convert the frequency; these are Shin Shinano, Sakuma Dam, Minami-Fukumitsu, and the Higashi-Shimizu Frequency Converter.
Utility frequencies in North America in 1897
Utility frequencies in Europe to 1900
Even by the middle of the 20th century, utility frequencies were still not entirely standardized at the now-common 50 Hz or 60 Hz. In 1946, a reference manual for designers of radio equipment listed the following now obsolete frequencies as in use. Many of these regions also had 50-cycle, 60-cycle, or direct current supplies.
Frequencies in use in 1946 (as well as 50 Hz and 60 Hz)
Where regions are marked (*), this is the only utility frequency shown for that region.
Railways
Other power frequencies are still used. Germany, Austria, Switzerland, Sweden, and Norway use traction power networks for railways, distributing single-phase AC at Hz or 16.7 Hz. A frequency of 25 Hz is used for the Austrian Mariazell Railway, as well as Amtrak and SEPTA's traction power systems in the United States. Other AC railway systems are energized at the local commercial power frequency, 50 Hz or 60 Hz.
Traction power may be derived from commercial power supplies by frequency converters, or in some cases may be produced by dedicated traction powerstations. In the 19th century, frequencies as low as 8 Hz were contemplated for operation of electric railways with commutator motors.
Some outlets in trains carry the correct voltage, but using the original train network frequency like Hz or 16.7 Hz.
400 Hz
Power frequencies as high as 400 Hz are used in aircraft, spacecraft, submarines, server rooms for computer power, military equipment, and hand-held machine tools. Such high frequencies cannot be economically transmitted long distances; the increased frequency greatly increases series impedance due to the inductance of transmission lines, making power transmission difficult. Consequently, 400 Hz power systems are usually confined to a building or vehicle.
Transformers, for example, can be made smaller because the magnetic core can be much smaller for the same power level. Induction motors turn at a speed proportional to frequency, so a high-frequency power supply allows more power to be obtained for the same motor volume and mass. Transformers and motors for 400 Hz are much smaller and lighter than at 50 or 60 Hz, which is an advantage in aircraft and ships. A United States military standard MIL-STD-704 exists for aircraft use of 400 Hz power.
Stability
Time error correction (TEC)
Regulation of power system frequency for timekeeping accuracy was not commonplace until after 1916 with Henry Warren's invention of the Warren Power Station Master Clock and self-starting synchronous motor. Nikola Tesla demonstrated the concept of clocks synchronized by line frequency at the 1893 Chicago Worlds fair. The Hammond Organ also depends on a synchronous AC clock motor to maintain the correct speed of its internal "tone wheel" generator, thus keeping all notes pitch-perfect.
Today, AC power network operators regulate the daily average frequency so that clocks stay within a few seconds of the correct time. In practice the nominal frequency is raised or lowered by a specific percentage to maintain synchronization. Over the course of a day, the average frequency is maintained at a nominal value within a few hundred parts per million. In the synchronous grid of Continental Europe, the deviation between network phase time and UTC (based on International Atomic Time) is calculated at 08:00 each day in a control center in Switzerland. The target frequency is then adjusted by up to ±0.01 Hz (±0.02%) from 50 Hz as needed, to ensure a long-term frequency average of exactly 50 Hz × 60 s/min × 60 min/h × 24 h/d = cycles per day. In North America, whenever the error exceeds 10 seconds for the Eastern Interconnection, 3 seconds for the Texas Interconnection, or 2 seconds for the Western Interconnection, a correction of ±0.02 Hz (0.033%) is applied. Time error corrections start and end either on the hour or on the half-hour.
Real-time frequency meters for power generation in the United Kingdom are available online – an official one for the National Grid, and an unofficial one maintained by Dynamic Demand. Real-time frequency data of the synchronous grid of Continental Europe is available on websites such as . The Frequency Monitoring Network (FNET) at the University of Tennessee measures the frequency of the interconnections within the North American power grid, as well as in several other parts of the world. These measurements are displayed on the FNET website.
US regulations
In the United States, the Federal Energy Regulatory Commission made time error correction mandatory in 2009. In 2011, The North American Electric Reliability Corporation (NERC) discussed a proposed experiment that would relax frequency regulation requirements for electrical grids which would reduce the long-term accuracy of clocks and other devices that use the 60 Hz grid frequency as a time base.
Frequency and load
Modern alternating-current grids use precise frequency control as an out-of-band signal to coordinate generators connected the network. The practice arose because the frequency of a mechanical generator varies with the input force and output load experienced. Excess load withdraws rotational energy from the generator shaft, reducing the frequency of the generated current; excess force deposits rotational energy, increasing frequency. Automatic generation control (AGC) maintains scheduled frequency and interchange power flows by adjusting the generator governor to counteract frequency changes, typically within several decaseconds.
Flywheel physics does not apply to inverter-connected solar farms or other DC-linked power supplies. However, such power plants or storage systems can be programmed to follow the frequency signal. Indeed, a 2017 trial for CAISO discovered that solar plants could respond to the signal faster than traditional generators, because they did not need to accelerate a rotating mass.
Small, temporary frequency changes are an unavoidable consequence of changing demand, but dramatic, rapid frequency shifts often signal that a distribution network is near capacity limits. Exceptional examples have occurred before major outages. During a severe failure of generators or transmission lines, the ensuing load-generation imbalance will induce variation in local power system frequencies. Loss of an interconnection causes system frequency to increase (due to excess generation) upstream of the loss, but may cause a collapse in frequency or voltage (due to excess load) downstream of the loss. Consequently many power system protective relays automatically trigger on severe underfrequency (typically too low, depending on the system's disturbance tolerance and the severity of protection measures). These initiate load shedding or trip interconnection lines to preserve the operation of at least part of the network.
Smaller power systems, not extensively interconnected with many generators and loads, will not maintain frequency with the same degree of accuracy. Where system frequency is not tightly regulated during heavy load periods, system operators may allow system frequency to rise during periods of light load to maintain a daily average frequency of acceptable accuracy. Portable generators, not connected to a utility system, need not tightly regulate their frequency because typical loads are insensitive to small frequency deviations.
Load-frequency control
Load-frequency control (LFC) is a type of integral control that restores the system frequency while respecting contracts for power provision or consumption to surrounding areas. The automatic generation scheme described in establishes a damping that minimizes the magnitude of average frequency error, , where is frequency, refers to the difference between measured and desired values, and overlines indicate time averages.
LFC incorporates power transfer between different areas, known as "net tie-line power", into the minimized quantity. For a particular frequency bias constant , the area control error (ACE) associated with LFC at any moment in time is simply where refers to tie-line power. This instantaneous error is then numerically integrated to give the time average, and governors adjusted to counteract its value. The coefficient traditionally has a negative value, so that when the frequency is lower than the target, area power production should increase; its magnitude is usually on the order of MW/dHz.
Tie-line bias LFC was known since 1930s, but was rarely used until the post-war period. In the 1950s, Nathan Cohn popularized the practice in a series of articles, arguing that load-frequency control minimized the adjustment necessary for changes in load. In particular, Cohn supposed that all regions of the grid shared a common linear regime, with location-invariant frequency change per additional loading (). If the utility selected and one region experienced a temporary fault or other generation-load mismatch, then adjacent generators would observe a decrease in frequency but a counterbalancing increase in outward tieline power flow, giving no ACE. They would thus make no governor adjustments in the (presumed) brief period before the failed region recovered.
Rate of change of frequency
Rate of change of frequency (also RoCoF) is simply a time derivative of the utility frequency (), usually measured in Hz per second, Hz/s. The importance of this parameter increases when the traditional synchronous generators are replaced by the variable renewable energy (VRE) inverter-based resources (IBR). The design of a synchronous generator inherently provides the inertial response that limits the RoCoF. Since the IBRs are not electromechanically coupled into the power grid, a system with high VRE penetration might exhibit large RoCoF values that can cause problems with the operation of the system due to stress placed onto the remaining synchronous generators, triggering of the protection devices and load shedding.
As of 2017, regulations for some grids required the power plants to tolerate RoCoF of 1–4 Hz/s, the upper limit being a very high value, an order of magnitude higher than the design target of a typical older gas turbine generator. Testing high-power (multiple MW) equipment for RoCoF tolerance is hard, as a typical test setup is powered off the grid, and the frequency thus cannot be arbitrarily varied. In the US, the controllable grid interface at the National Renewable Energy Laboratory is the only facility that allows testing of multi-MW units (up to 7 MVA). Testing of large thermal units is not possible.
Audible noise and interference
AC-powered appliances can give off a characteristic hum, often called "mains hum", at the multiples of the frequencies of AC power that they use (see Magnetostriction). It is usually produced by motor and transformer core laminations vibrating in time with the magnetic field. This hum can also appear in audio systems, where the power supply filter or signal shielding of an amplifier is not adequate.
Most countries chose their television vertical synchronization rate to be the same as the local mains supply frequency. This helped to prevent power line hum and magnetic interference from causing visible beat frequencies in the displayed picture of early analogue TV receivers particularly from the mains transformer. Although some distortion of the picture was present, it went mostly un-noticed because it was stationary. The elimination of transformers by the use of AC/DC receivers, and other changes to set design helped minimise the effect and some countries now use a vertical rate that is an approximation to the supply frequency (most notably 60 Hz areas).
Another use of this side effect is as a forensic tool. When a recording is made that captures audio near an AC appliance or socket, the hum is also incidentally recorded. The peaks of the hum repeat every AC cycle (every ms for 50 Hz AC, or every ms for 60 Hz AC). The exact frequency of the hum should match the frequency of a forensic recording of the hum at the exact date and time that the recording is alleged to have been made. Discontinuities in the frequency match or no match at all will betray the authenticity of the recording.
| Technology | Concepts | null |
488945 | https://en.wikipedia.org/wiki/Zebra%20mussel | Zebra mussel | The zebra mussel (Dreissena polymorpha) is a small freshwater mussel, an aquatic bivalve mollusk in the family Dreissenidae. The species originates from the lakes of southern Russia and Ukraine, but has been accidentally introduced to numerous other areas and has become an invasive species in many countries worldwide. Since the 1980s, the species has invaded the Great Lakes, Hudson River, Lake Travis, Finger Lakes, Lake Bonaparte, and Lake Simcoe. The adverse effects of dreissenid mussels on freshwater systems have led to their ranking as one of the world's most invasive aquatic species.
The species was first described in 1769 by German zoologist Peter Simon Pallas in the Ural, Volga, and Dnieper Rivers. Zebra mussels get their name from a striped pattern commonly seen on their shells, though it is not universally present. They are usually about the size of a fingernail, but can grow to a maximum length around . Their shells are D-shaped, and attached to the substrate with strong byssal fibers, which come out of their umbo on the dorsal (hinged) side.
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Right and left valve of the same specimen:
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Ecology
Zebra mussels and the closely related and ecologically similar quagga mussels are filter-feeding organisms; they remove particles from the water column. Zebra mussels process up to of water per day, per mussel. Some particles are consumed as food, and feces are deposited on the lake floor. Nonfood particles are combined with mucus and other matter and deposited on lake floors as pseudofeces. Since the zebra mussel has become established in Lake Erie, water clarity has increased from to up to in some areas. This increased water clarity allows sunlight to penetrate deeper, enabling growth of submerged macrophytes. These plants, when decaying, wash up on shorelines, fouling beaches and causing water-quality problems.
Lake floor food supplies are enriched by zebra mussels as they filter pollution out of the water. This biomass becomes available to bottom-feeding species and to the fish that feed on them. The catch of yellow perch increased five-fold after the invasion of zebra mussels into Lake St. Clair.
Zebra mussels attach to most substrates, including sand, silt, and harder substrates, but usually juveniles prefer harder, rockier substrates on which to attach. Other mussel species frequently represent the most stable objects in silty substrates, and zebra mussels attach to and often kill these mussels. They build colonies on native unionid clams, reducing their ability to move, feed, and breed, eventually leading to their deaths. This has led to the near extinction of the unionid clams in Lake St. Clair and the western basin of Lake Erie. This pattern is being repeated in Ireland, where zebra mussels have eliminated the two freshwater mussels from several waterways, including some lakes along the River Shannon in 1997.
In 2012, the National University of Ireland, Galway, said "the discovery of zebra mussels (Dreissena polymorpha) in Lough Derg and the lower Shannon region in 1997 has led to considerable concern about the potential ecological and economic damage that this highly invasive aquatic nuisance species can cause."
The zebra mussel is a freshwater species and cannot survive in the ocean.
Life cycle
The lifespan of a zebra mussel is four to five years. A female zebra mussel begins to reproduce within 6–7 weeks of settling.
An adult female zebra mussel can produce 30,000 to 40,000 eggs in each reproductive cycle, and over 1 million each year. Free-swimming microscopic larvae, called veligers, drift in the water for several weeks and then settle onto any hard surface they can find. Zebra mussels also can tolerate a wide range of environmental conditions, and adults can even survive out of water for about seven days.
Indicators of environmental pollution
Zebra mussels can be used to detect risks to humans from environmental hazards and are considered effective indicators of environmental pollution. Exposure of D. polymorpha gill cells to model contaminants representative of water pollution was found to cause DNA damage measured as the formation of 8-oxodG and DNA strand breaks.
Predators
Research on natural enemies, both in Europe and North America, has focused on predators, particularly birds (36 species) and fish (15 species eating veligers and 38 eating attached mussels). Annually, the wintering waterbirds at Lake Constance decrease zebra mussel biomass in shallow areas by >90%. Biomass reduction in deeper areas varies considerably based on substratum; Werner et al. observed no reduction at the lowest observed depth of except for a site at Hagnau.
The vast majority of the organism's natural enemies are not present in North America. Ecologically similar species do exist, but these species are unlikely to be able to eliminate those mussels already established and have a limited role in their control unlike their counterparts in Europe.
It is pointed out that crayfish could have a significant impact on the densities of -long zebra mussels. An adult crayfish consumes around 105 zebra mussels every day, or about 6,000 mussels in a season. However, predation rates are significantly reduced at lower water temperatures. Additionally, certain fish, such as the Smallmouth bass, is a predator in the zebra mussels' adopted North American Great Lakes habitat, but in European lakes, fish do not seem to limit the densities of zebra mussels.
There have been some high winter mortalities, for example in the winter of 1994–1995 in the invasive population of Lake Simcoe. Evans et al., 2011 attributes this to predation by the crayfish Orconectes propinquus. Nonetheless this has not been sufficient to eradicate the problem.
Other control
On June 4, 2014, Canadian conservation authorities announced that a test using liquid fertilizer to kill invasive zebra mussels was successful. This test was conducted in a lakefront harbor in the western province of Manitoba. However, outbreaks continue in Lake Winnipeg.
Similar tests were run in Illinois, Minnesota, and Michigan, using zequanox, a biopesticide. Niclosamide proves effective in killing invasive zebra mussels in cool waters.
As an invasive species
Europe
The native distribution of the species is in the Black Sea and Caspian Sea in Eurasia. Zebra mussels have become an invasive species in North America, Great Britain, Ireland, Italy, Spain, and Sweden. They disrupt the ecosystems by monotypic colonization, and damage harbors and waterways, ships and boats, and water-treatment and power plants. Water-treatment plants are most affected because the water intakes bring the microscopic, free-swimming larvae directly into the facilities. Zebra mussels also cling to pipes under the water and clog them.
Grossinger reported it in Hungary in 1794. Kerney and Morton described the rapid colonization of Britain by the zebra mussel, first in Cambridgeshire in the 1820s, London in 1824, and in the Union Canal near Edinburgh in 1834.
In 1827, zebra mussels were seen in the Netherlands at Rotterdam. Canals that artificially link many European waterways facilitated their early dispersal. It is nonindigenous in the Czech Republic in the Elbe River in Bohemia since 1893; in southern Moravia, it is probably native. Around 1920 the mussels reached Lake Mälaren in Sweden.
The first appearance of the organism in northern Italy was in Lake Garda in 1973; in central Italy, they appeared in Tuscany in 2003.
Zebra mussels are present in British waterways. Many water companies are reporting having problems with their water-treatment plants with the mussels attaching themselves to pipeworks. Anglian Water has estimated that it costs £500,000 per year to remove the mussels from their treatment plants. Zebra mussels arguably have also had an effect on fishing, for example at Salford Quays, where their introduction has changed the environment for the fish.
Zebra mussels were first reported in Ireland in 1997, but probably arrived in 1994 or earlier. First identified near Lough Derg, they may have been introduced to Ireland through several vectors but have since spread through much of the River Shannon and its adjacent waters. In the summers of 2023 and 2024, zebra mussels were one contributing factor to a major bloom of toxic cyanobacteria in Lough Neagh, the largest body of freshwater in the United Kingdom. In 2021 the similar quagga mussel was identified in the Shannon, which is capable of tolerating a wider range of environmental conditions. The spread of quagga mussels is expected to reduce zebra mussel populations over time while increasing the overall environmental impact.
The mussels have displaced native species of molluscs in Lake Constance, reaching densities of up to . The mussels present a food source to waterfowl and have caused bird numbers to double over the last 30 years. By the end of winter, birds decimate zebra mussel populations and reduce them by 95–99% up to the maximum depth reachable by birds of c. . The estimated quantity of consumed zebra mussels is . Zebra mussel populations recover annually, indicating that waterfowl may control infested bodies of water but not reverse the infestation status entirely.
North America
They were first detected in Canada in the Great Lakes in 1988, in Lake St. Clair. They are thought to have been inadvertently introduced into the lakes by the ballast water of ocean-going ships that were traversing the St. Lawrence Seaway. Another possible, but unproven, mode of introduction is on anchors and chains. Since adult zebra mussels can survive out of water for several days or weeks if the temperature is low and humidity is high, chain lockers provide temporary refuge for clusters of adult mussels that could easily be released when transoceanic ships drop anchor in freshwater ports. They have become an invasive species in North America, and as such, they are the target of federal policy to control them, for instance in the National Invasive Species Act (1996).
Using models based on the genetic algorithm for rule-set production (GARP), a group of researchers predicted that the Southeastern United States is moderately to highly likely to be inhabited by zebra mussels and the Midwest unlikely to experience a zebra mussel invasion of water bodies. This model, though, has since been proven incorrect. In 2006, a researcher (also using GARP) predicted invasion as far west as the North Platte River by 2015. As of March 2016, zebra mussels have affected hundreds of lakes in the Midwest including Lake Michigan, and the largest interior lake in Wisconsin, Lake Winnebago.
Congressional researchers have estimated that the zebra mussel has cost businesses and communities over $5 billion since their initial invasion. Zebra mussels have cost power companies alone over $3 billion.
On 2 March 2021, the US Geological Survey was notified that zebra mussels had been discovered in marimo moss balls, a common aquarium plant, sold in pet stores across North America. By 8 March, the invasive species were detected in moss balls in 30 different states at multiple retail locations in the United States. These discoveries were prompted by the initial find at a Petco in Seattle. Infested-moss balls have also been found from online retailers and smaller, independent stores. Though it is more difficult to know the extent of the spread at a larger scale, Wesley Daniel, a fisheries biologist with the U.S. Geological Survey, says that about 30% of the inventory pulled from shelves were found to contain the zebra mussels.
After working with the USGS, PetSmart and Petco voluntarily recalled their moss balls due to the potential harm zebra mussels could cause to indigenous ecosystems. As of November 2021, nearly eight months after the recall, marimo moss balls have not returned to shelves.
By location
From their first appearance in North American waters in 1988, zebra mussels have spread to a large number of waterways, including Lake Simcoe in the Great Lakes region, the Mississippi, Hudson, St. Lawrence, Ohio, Cumberland, Missouri, Tennessee, Huron, Colorado, and Arkansas rivers, and, by mid-2023, 31 lakes and seven river basins in Texas. However, zebra mussels could not establish reproducing populations in Kentucky Lake (Tennessee River). This may be due to relatively low dissolved calcium levels in this water body.
In 2009, the Massachusetts Department of Conservation and Recreation confirmed that zebra mussels had been found in Laurel Lake in the Berkshires. That same year, the Minnesota Department of Natural Resources announced that live zebra mussels had been found in Pelican Lake. This was the first confirmed sighting in the Red River Basin, which extends across the international border into the province of Manitoba. In 2013, their presence in Manitoba's Lake Winnipeg was confirmed, and aggressive efforts to eradicate them in 2014 have not succeeded. New contamination was found outside treated areas of Lake Winnipeg in 2015, and they have also been found in the Red River near the lake in Selkirk Park in 2015. Large numbers were seen at Grand Beach in 2017.
In July 2010, the North Dakota Game and Fish Department confirmed the presence of zebra mussel veliger in the Red River between Wahpeton, North Dakota, and Breckenridge, Minnesota. , California similarly reported invasions.
In 2011, an invasion of zebra mussels shut down a water pipeline in the Dallas area. This resulted in reduced water supplies during a drought year, worsening water restrictions across the Dallas area.
A common inference made by scientists predicts that the zebra mussel will continue spreading passively, by ship and by pleasure craft, to more rivers in North America. Trailered boat traffic is the most likely vector for invasion into Western North America. This spread is preventable if boaters thoroughly clean and dry their boats and associated equipment before transporting them to new bodies of water. Since no North American predator or combination of predators has been shown to significantly reduce zebra mussel numbers, such spread would most likely result in permanent establishment of zebra mussels in many North American waterways.
A major decrease in the concentration of dissolved oxygen was observed in the Seneca River in central New York in the summer of 1993. This decrease was caused by extremely high concentrations of zebra mussels in the watershed. Additionally, the Seneca River had significantly less chlorophyll in the water, which is used as a measure of phytoplankton biomass, because of the presence of zebra mussels.
Cost
The cost of fighting the pests at power plants and other water-consuming facilities is substantial, but the magnitude of the damages is a matter of some controversy. According to the Center for Invasive Species Research at the University of California, Riverside, the cost of management of zebra mussels in the Great Lakes alone exceeds $500 million per year. A more conservative study estimated total economic costs of $267 million for electric-generation and water-treatment facilities in the entire United States from 1989 through 2004. In a study conducted by the US Department of State in 2009, the total cost of the zebra mussel invasion is estimated at 3.1 billion over the next 10 years.
Concerns are also high following the contamination of zebra mussels in at home aquariums. If zebra mussels had reached open water in Seattle, Washington, where the first case was confirmed, the invasive species would have cost the state $100 million each year in maintenance for power and water systems.
Effects
As with most bivalves, zebra mussels are filter feeders. When in the water, they open their shells to admit detritus. Their shells are very sharp and are known for cutting people's feet, resulting in the need to wear water shoes wherever they are prevalent.
Since their colonization of the Great Lakes, they have covered the undersides of docks, boats, and anchors. They have also spread into streams and rivers throughout the U.S. In some areas, they completely cover the substrate, sometimes covering other freshwater mussels. They can grow so densely that they block pipelines, clogging water intakes of municipal water supplies and hydroelectric companies. Zebra mussels do not attach to cupronickel alloys, which can be used to coat intake and discharge grates, navigational buoys, boats, and motors where the species tends to congregate.
Zebra mussels are believed to be the source of deadly avian botulism poisoning that has killed tens of thousands of birds in the Great Lakes since the late 1990s. They are edible, but since they are so efficient at filtering water, they tend to accumulate pollutants and toxins, so most experts recommend against consuming zebra mussels.
Zebra mussels affect all classes of algal species, resulting in a shortage of food sources to native species of freshwater mussels and fish in the Great Lakes.
However, zebra mussels and other non-native species are credited with the increased population and size of smallmouth bass in Lake Erie and yellow perch in Lake St. Clair. They cleanse the waters of inland lakes, resulting in increased sunlight penetration and growth of native algae at greater depths. This cleansing also increases water visibility and filters out pollutants. Each quagga and zebra mussel filters about of water per day when confined to small tanks. In lakes, their filtering effects are usually spatially restricted (near the lake bottom) because of nonhomogeneous water column mixing.
In Lake Pleshcheyevo in Russia, zebra mussels have greatly changed the fish community. All fish species remained in the lake, but catches changed significantly. Gillnets set in the littoral and sublittoral zones during the feeding period consist mainly of large roach and perch while vendace were prevalent in the pelagic zone. The abundance of benthophagous fishes increased slightly due to the presence of the zebra mussel which made up a significant part of the diet followed by an increase (p<0.05) in the growth rates of roach and silver bream. The growth rates of the bream, which prefers soft zoobenthos, decreased. In addition, the formation of a stable biocenosis of zebra mussels in the lake probably caused changes in the spatial structure of the fish community. What is notable is the disappearance of small roach from pelagic assemblages, probably due to alterations in the trophic links of the littoral and sublittoral zones, as well as the expected increase in food competition among pelagic fish species.
Because zebra mussels damage water intakes and other infrastructure, methods such as adding oxidants, flocculants, heat, dewatering, mechanical removal, and pipe coatings are becoming increasingly common.
Preventing their spread
Zebra mussels cling to boat motors. Boat-owners should follow a few steps prior to putting their boats into a new lake and after removing their boats from infected lakes to stop the spreading of the species. Boat owners should make sure to inspect their boat, trailer, and other recreational equipment that have been in contact with water, remove all mud, plants, or animals, drain all bilge water, live wells, bait buckets, and all other water from their boats, engines and equipment, wash all parts of their boats, paddles, and other equipment that have been in contact with water, and dry their boats and trailers in the sun for five days before launching into another body of water. This is important because adult zebra mussels are able to close their shells and may survive out of water for several days. When washing their boats, boat owners should be sure to wash with warm, soapy water as well.
If marimo moss balls were purchased around the time of the first discovery and recall in March 2021, aquatic hobbyists have been urged to decontaminate the moss balls by either boiling them for at least one minute, freezing for at least 24 hours, or placing them in diluted chlorine bleach. Another way to rid of the zebra mussels can be to submerge the moss balls in undiluted white vinegar for a minimum of 20 minutes. After following one of these methods, the USGS urges owners to bag these moss balls before disposing of them in the trash to prevent spread to local water ways and ecosystems.
Even if moss balls infected with zebra mussels are contained in an aquarium, the concern that they could contaminate local waterways is high, especially in regions and states where they have not yet infested. Aquarium dumping and disposing of unwanted pets is common, according to Eric Fischer with Indiana's Department of Natural Resources (DNR). It is illegal to own, sell, or distribute zebra mussels in the U.S. If spotted, either in an aquarium or out in nature, it is advised to contact the local DNR in the region resided. Recent research shows that zebra mussels could not establish reproducing populations in Kentucky Lake, indicating that the physico-chemical characteristics of this lake deter the species. A better understanding of the lake's characteristics and identification of the key parameters that deter zebra mussels may pave the way for protecting other aquatic ecosystems from the spreading of this invasive species.
| Biology and health sciences | Bivalvia | Animals |
489079 | https://en.wikipedia.org/wiki/Swiftlet | Swiftlet | Swiftlets are birds from the four genera Aerodramus, Collocalia, Hydrochous and Schoutedenapus, which form the tribe Collocaliini within the swift family Apodidae. The group contains around thirty species mostly confined to southern Asia, south Pacific islands, and northeastern Australia, all within the tropical and subtropical regions. They are in many respects typical members of the Apodidae, having narrow wings for fast flight, with a wide gape and small reduced beak surrounded by bristles for catching insects in flight. What distinguishes many but not all species from other swifts and indeed almost all other birds is their ability to use a simple but effective form of echolocation to navigate in total darkness through the chasms and shafts of the caves where they roost at night and breed.
The nests of some species are built entirely from solidified threads of their saliva, which are edible and thus collected for human consumption as the famous Chinese delicacy, the bird's nest soup.
Description and ecology
The swift family remains one of the more complicated groups of birds in taxonomic research, but the swiftlet tribe is a rather well-defined group. Its internal systematics is confusing; the plumage is usually dull, with shades of black, brown, and gray; from their outward appearance, most species are very similar. Swiftlets have four toes, except the Papuan swiftlet which lacks the hallux (back toe). Their legs are very short, preventing the birds from perching, but allowing them to cling to vertical surfaces. Flight is mainly gliding due to very long primary feathers and small breast muscles. The larger Aerodramus swiftlets weigh about 14 grams and are 10 cm long.
Swiftlets are insectivores; hymenopterans and dipterans being the most abundant prey. Typically, they leave the cave during the day to forage and return to their roost at night. Males and females look similar; as usual in such cases, these birds are monogamous and both partners take part in caring for the nestlings. Males perform aerial displays to attract females and mating occurs at the nest. The breeding season overlaps the wet season, which corresponds to an increased insect population. Clutch size depends on the location and the food source, but it is generally not large; Aerodramus swiftlets lay 1 to 2 eggs. The eggs are a dull white color and are laid every other day. Many if not all species are colonial nesters; some build their nests in high, dark corners on cave walls. Swiftlets in temperate zones do migrate, but most Aerodramus swiftlets live in the tropical Indo-Pacific region and do not migrate. These birds usually remain in one cave or other roosting/nesting site. Some examples of caves include the Niah Caves at Niah National Park & Gunung Mulu National Park which are all located in Sarawak, Malaysian Borneo.
The genus Aerodramus is of special interest due to its use of echolocation and its intricately constructed saliva nests which in some species contain no other material such as feathers, moss or twigs and are collected, selling at extremely high prices (see Bird's nest soup). It has been argued that the high demand for these nests could have had an adverse effect on their populations, but other authorities have shown that modern techniques of nest farming have increased the bird population.
The use of echolocation was once used to separate Aerodramus from the non-echolocating genera Collocalia and Hydrochous (virtually nothing is known about Schoutedenapus). But recently, the pygmy swiftlet Collocalia troglodytes was discovered making similar clicking noises in and outside its cave. Characteristics of behavior, such as what materials other than saliva the nests contain, can be used to differentiate between certain species of Aerodramus.
Echolocation
The genus Aerodramus was thought to be the only echolocating swiftlets. These birds use echolocation to locate their roost in dark caves. Unlike a bat's echolocation, Aerodramus swiftlets make clicking noises that are well within the human range of hearing. The clicks consist of two broad band pulses (3–10 kHz) separated by a slight pause (1–3 milliseconds). The interpulse periods (IPPs) are varied depending on the level of light; in darker situations the bird emits shorter IPPs, as obstacles become harder to see, and longer IPPs are observed when the bird nears the exit of the cave. This behavior is similar to that of bats as they approach targets. The birds also emit a series of low clicks followed by a call when approaching the nests; presumably to warn nearby birds out of their way. It is thought that the double clicks are used to discriminate between individual birds. Aerodramus sawtelli, the Atiu swiftlet, and Aerodramus maximus, the black-nest swiftlet are the only known species which emit single clicks. The single click is thought be used to avoid voice overlap during echolocation. The use of a single click might be associated with an evolutionary shift in eastern Pacific swiftlets; determining how many clicks the Marquesan swiftlet emits could shed light on this. It was also discovered that both the Atiu swiftlet and the Papuan swiftlet emit clicks while foraging outside at dusk; the latter possibly only in these circumstances, considering that it might not nest in caves at all. Such behavior is not known to occur in other species, but quite possibly does, given that the Papuan and Atiu swiftlets are not closely related. However, it has recently been determined that the echolocation vocalizations do not agree with evolutionary relationship between swiftlet species as suggested by DNA sequence comparison. This suggests that as in bats, echolocation sounds, once present, adapt rapidly and independently to the particular species' acoustic environment.
Three hypotheses are considered to describe how echolocation evolved in the genus Aerodramus and, as determined more recently, other taxa in the Apodidae. One hypothesis states that echolocation evolved from an ancestral species of swiftlets and was lost in the genera which lack echolocation. A second hypothesis is that echolocation evolved independently several times. The third scenario involves a combination of the first two, i.e. a gain-loss-regain scenario.
Several functional subunits (like vocal muscles and brain areas) are needed to produce the echolocating system. Past studies have thought that the loss of one of these subunits was more likely to occur than acquiring all the traits needed to echolocate. Yet a recent study suggests that the echolocation subunits were mainly located in the central nervous system, while the subunits in the vocal apparatus were already present and capable of use before echolocation even evolved. This study supports the second hypothesis of independent evolution of echolocation in Aerodramus and Collocalia, with the subsequent evolution of complex behavior needed to complement the physical echolocation system, or even the third approach, as the vocal apparatus-parts of the echolocation system might even be inherited from some prehistoric nocturnal ancestor.
Culinary use
Authentic bird's-nest soup is made from nests of some species of swiftlet, mainly the edible-nest (or white-nest) swiftlet (Aerodramus fuciphagus) and the black-nest swiftlet. Instead of twigs, feathers and straw, these swiftlets make their nest only from strands of their gummy saliva, which hardens when exposed to air. Once the nests are harvested, they are cleaned and sold to restaurants. Eating swiftlet nest material is believed to help maintain skin tone, balance qi ("life energy") and reinforce the immune system,. It is also believed to strengthen the lungs and prevent coughs, improve the constitution and prolong life. The nutritional value of 100 g of dry nest includes 49.9 g of water-soluble protein (including amido nitrogen, monoamine nitrogen, non-amino nitrogen, arginine, humin, histidine, lysine and cysteine), 30.6 g carbohydrate (glycoprotein and mucin), 4.9 g iron, 2.5 g inorganic salt (including potassium, sodium, calcium, magnesium, sulfur, phosphorus, silica and other trace elements), and 1.4 g fiber (Dictionary of Traditional Chinese Medicine, The History of Chinese Medicine and the Nutrition Table).
The energy contained in 100 g of swiftlet nest is 345 kcal. The nests are often served simmered in chicken broth.
Authentic bird's-nest soup is quite popular throughout Asia. It is also extremely expensive; many western restaurants serve a less expensive version consisting of soup with noodles shaped to resemble a bird's nest.
Cave ecology
Guano from both the swiftlets and the many bats that inhabit the caves supports a huge array of specialized animals that feed on the dung. There are yet other creatures that have evolved to feed on these dung eaters as well as on the bats and the swiftlets themselves, including snakes that can climb the sheer walls to snatch a passing meal and huge carnivorous crickets that prey on chicks and bat pups. This cave fauna ecosystem is self-sustaining, the only link with the outside being the birds and the bats that bring the nutrients into the caves in the first place.
The Philippine municipality of El Nido in Palawan, known for its limestone cliffs and pristine beaches, is home to a thriving bird's-nest market. The name El Nido is the Spanish term for literally "The Nest". Many locals still practice manual climbing of the limestone caves to gather swiftlet nests.
Species
The Papuan swiftlet is apparently closer to the waterfall swift than to the other Aerodramus species and probably best placed in a separate genus, whereas Thomassen et al. (2005) advocate reuniting all swiftlets in Collocalia. Schoutedenapus is one of the least-known genera of birds.
Genus Collocalia
Plume-toed swiftlet, Collocalia affinis
Grey-rumped swiftlet, Collocalia marginata
Ridgetop swiftlet, Collocalia isonota
Tenggara swiftlet, Collocalia sumbawae
Drab swiftlet, Collocalia neglecta
Glossy swiftlet, Collocalia esculenta
Satin swiftlet, Collocalia uropygialis
Bornean swiftlet, Collocalia dodgei
Cave swiftlet, Collocalia linchi
Christmas Island swiftlet, Collocalia natalis
Pygmy swiftlet, Collocalia troglodytes
Genus Aerodramus
Seychelles swiftlet, Aerodramus elaphrus
Mascarene swiftlet, Aerodramus francicus
Indian swiftlet, Aerodramus unicolor
Philippine swiftlet, Aerodramus mearnsi
Moluccan swiftlet, Aerodramus infuscatus
Mountain swiftlet, Aerodramus hirundinaceus
White-rumped swiftlet, Aerodramus spodiopygius
Australian swiftlet, Aerodramus terraereginae
Himalayan swiftlet, Aerodramus brevirostris
Indochinese swiftlet, Aerodramus rogersi
Volcano swiftlet, Aerodramus vulcanorum
Whitehead's swiftlet, Aerodramus whiteheadi
Bare-legged swiftlet, Aerodramus nuditarsus
Mayr's swiftlet, Aerodramus orientalis
Palawan swiftlet, Aerodramus palawanensis
Mossy-nest swiftlet, Aerodramus salangana
Uniform swiftlet, Aerodramus vanikorensis
Palau swiftlet, Aerodramus pelewensis
Guam swiftlet, Aerodramus bartschi
Caroline Islands swiftlet, Aerodramus inquietus
Mangaia swiftlet, Aerodramus manuoi (prehistoric)
Atiu swiftlet, Aerodramus sawtelli
Polynesian swiftlet, Aerodramus leucophaeus
Marquesan swiftlet, Aerodramus ocistus
Black-nest swiftlet, Aerodramus maximus
Edible-nest swiftlet, Aerodramus fuciphagus
Brown-rumped swiftlet, Aerodramus (fuciphagus) vestitus
German's swiftlet, Aerodramus germani
Papuan swiftlet, Aerodramus papuensis – probably a distinct genus
Genus Hydrochous
Giant swiftlet, Hydrochous gigas
Genus Schoutedenapus
Scarce swift, Schoutedenapus myoptilus
| Biology and health sciences | Apodiformes | Animals |
489173 | https://en.wikipedia.org/wiki/Peripheral%20artery%20disease | Peripheral artery disease | Peripheral artery disease (PAD) is a vascular disorder that causes abnormal narrowing of arteries other than those that supply the heart or brain. PAD can happen in any blood vessel, but it is more common in the legs than the arms.
When narrowing occurs in the heart, it is called coronary artery disease (CAD), and in the brain, it is called cerebrovascular disease. Peripheral artery disease most commonly affects the legs, but other arteries may also be involved, such as those of the arms, neck, or kidneys.
Peripheral artery disease (PAD) is a form of peripheral vascular disease. Vascular refers to both the arteries and veins within the body. PAD differs from peripheral veinous disease. PAD means the arteries are narrowed or blocked—the vessels that carry oxygen-rich blood as it moves away from the heart to other parts of the body. Peripheral veinous disease, on the other hand, refers to problems with veins—the vessels that bring the blood back to the heart.
The classic symptom is leg pain when walking, which resolves with rest and is known as intermittent claudication. Other symptoms include skin ulcers, bluish skin, cold skin, or abnormal nail and hair growth in the affected leg. Complications may include an infection or tissue death, which may require amputation; coronary artery disease; or stroke. Up to 50% of people with PAD do not have symptoms.
The greatest risk factor for PAD is cigarette smoking. Other risk factors include diabetes, high blood pressure, kidney problems, and high blood cholesterol. PAD is primarily caused by the buildup of fatty plaque in the arteries, which is called atherosclerosis, especially in individuals over 40 years old. Other mechanisms include artery spasm, blood clots, trauma, fibromuscular dysplasia, and vasculitis. PAD is typically diagnosed by finding an ankle-brachial index (ABI) less than 0.90, which is the systolic blood pressure at the ankle divided by the systolic blood pressure of the arm. Duplex ultrasonography and angiography may also be used. Angiography is more accurate and allows for treatment at the same time; however, it is associated with greater risks.
It is unclear if screening for peripheral artery disease in people without symptoms is useful, as it has not been properly studied. For those with intermittent claudication from PAD, stopping smoking and supervised exercise therapy may improve outcomes. Medications, including statins, ACE inhibitors, and cilostazol, may also help. Aspirin, which helps with thinning the blood and thus improving blood flow, does not appear to help those with mild disease but is usually recommended for those with more significant disease due to the increased risk of heart attacks. Anticoagulants (blood thinners) such as warfarin show no definitive scientific evidence of benefit in PAD. Surgical procedures used to treat PAD include bypass grafting, angioplasty, and atherectomy.
In 2015, about 155 million people had PAD worldwide. It becomes more common with age. In the developed world, it affects about 5.3% of 45- to 50-year-olds and 18.6% of 85- to 90-year-olds. In the developing world, it affects 4.6% of people between the ages of 45 and 50 and 15% of people between the ages of 85 and 90. PAD in the developed world is equally common among men and women, though in the developing world, women are more commonly affected. In 2015, PAD resulted in about 52,500 deaths, which is an increase from the 16,000 deaths in 1990.
Signs and symptoms
The signs and symptoms of peripheral artery disease are based on the part of the body that is affected. About 66% of patients affected by PAD either do not have symptoms or have atypical symptoms. The most common presenting symptom is intermittent claudication (IC), which typically refers to lower extremity skeletal muscle pain that occurs during exercise. IC presents when there is insufficient oxygen delivery to meet the metabolic requirements of the skeletal muscles. IC is a common manifestation of peripheral arterial disease (PAD). The pain is usually located in the calf muscles of the affected leg and is relieved by rest. This occurs because during exercise, the muscles require more oxygen. Normally, the arteries would be able to increase the amount of blood flow and therefore increase the amount of oxygen going to the exercised muscle. However, in PAD, the artery is unable to respond appropriately to the increased demand for oxygen from the muscles, and as a result, the muscles are deprived of oxygen, leading to muscle pain that subsides with rest.
Other symptoms may include:
Pain, aches, and/or cramps in the buttocks, hip, or thigh
Muscle atrophy (muscle loss) of the affected limb
Hair loss of the affected limb
Skin that is smooth, shiny, or cool to the touch in the affected area
Decreased or absent pulse in the feet
Cold and/or numbness in the toes
Sores/ulcers on the affected limb that do not heal
In individuals with severe PAD, complications may arise, including critical limb ischemia and gangrene. Critical limb ischemia occurs when the obstruction of blood flow in the artery is compromised to the point where the blood is unable to maintain oxygenation of the tissue at rest. This can lead to pain at rest, a feeling of coldness, or numbness in the affected foot and toes. Other complications of severe PAD include lower limb tissue loss (amputation), arterial insufficiency ulcers, erectile dysfunction, and gangrene. People with diabetes are affected by gangrene of the feet at a rate that is 30 times higher than the unaffected population. Many of these severe complications, such as those leading to amputation, are irreversible.
Causes
Risk factors
Factors contributing to an increased risk of PAD are the same as those for atherosclerosis. These include age, sex, and ethnicity. PAD is twice as common in males as in females. In terms of ethnicity, PAD is more common in people of color compared to the white population in a 2:1 ratio. The factors with the greatest risk associations are hyperlipidemia, hypertension, diabetes mellitus, chronic kidney disease, and smoking. Presenting three of these factors or more increases the risk of developing PAD tenfold.
Smoking – Tobacco use in any form is the single greatest risk factor for peripheral artery disease internationally. Smokers have up to a 10-fold increase in the risk of PAD in a dose-response relationship. Exposure to second-hand smoke has also been shown to promote changes in the lining of blood vessels (endothelium), which can lead to atherosclerosis. Smokers are 2–3 times more likely to have lower extremity PAD than coronary artery disease. Greater than 80%–90% of patients with lower extremity peripheral arterial disease are current or former smokers. The risk of PAD increases with the number of cigarettes smoked per day and the number of years smoked.
High blood sugar – Diabetes mellitus is shown to increase the risk of PAD by 2–4 fold. It does this by causing endothelial and smooth-muscle cell dysfunction in peripheral arteries. The risk of developing lower extremity peripheral arterial disease is proportional to the severity and duration of diabetes.
High blood cholesterol – Dyslipidemia, which is an abnormally high level of cholesterol or fat in the blood. Dyslipidemia is caused by a high level of a protein called low-density lipoprotein (LDL cholesterol), low levels of high-density lipoprotein (HDL cholesterol), elevation of total cholesterol, and/or high triglyceride levels. This abnormality in blood cholesterol levels has been correlated with accelerated peripheral artery disease. Management of dyslipidemia by diet, exercise, and/or medication is associated with a major reduction in rates of heart attack and stroke.
High blood pressure – Hypertension or elevated blood pressure can increase a person's risk of developing PAD. Similarly to PAD, there is a known association between high blood pressure and heart attacks, strokes, and abdominal aortic aneurysms. High blood pressure increases the risk of intermittent claudication, the most common symptom of PAD, by 2.5- to 4-fold in men and women, respectively.
Other risk factors that are being studied include levels of various inflammatory mediators such as C-reactive protein, fibrinogen, homocysteine, and lipoprotein A. Individuals with increased levels of homocysteine in their blood have a 2-fold risk of developing peripheral artery disease. While there are genetic factors leading to risk factors for peripheral artery disease, including diabetes and high blood pressure, there have been no specific genes or gene mutations directly associated with the development of peripheral artery disease.
High risk populations
Peripheral arterial disease is more common in these populations:
All people who have leg symptoms with exertion (suggestive of claudication) or ischemic rest pain
All people aged 65 years and over, regardless of risk factor status
All people between 50 and 69 who have a cardiovascular risk factor (particularly diabetes or smoking)
Age less than 50 years, with diabetes and one other atherosclerosis risk factor (smoking, dyslipidemia, hypertension, or hyperhomocysteinemia)
Individuals with an abnormal lower extremity pulse examination
Those with known atherosclerotic coronary, carotid, or renal artery disease
All people with a Framingham risk score of 10%–20%
All people who have previously experienced chest pain
Etiology and pathophysiology
Peripheral arterial disease is considered to be a set of chronic or acute syndromes, generally derived from the presence of occlusive arterial disease, which causes inadequate blood flow to the limbs.
As previously mentioned, the most common etiology of peripheral artery disease, especially in patients over 40 years old, is atherosclerosis. Atherosclerosis is a narrowing of the arteries caused by lipid or fat buildup and calcium deposition in the wall of the affected arteries.
The pathophysiology of atherosclerosis involves complex interactions between cholesterol and vascular cells. In the early stages of PAD, the arteries compensate for the plaque buildup by dilating to preserve flow through the vessel. Eventually, the artery cannot dilate any further, and the atherosclerotic plaque starts to narrow the arterial flow lumen.
When there is an imbalance between the needs of the peripheral tissues and the blood supply, we are faced with a situation of ischemia.
From the pathophysiologic point of view, a restriction of blood supply (ischemia) to the lower limbs can be classified as either functional or critical. Functional ischemia occurs when the blood flow is normal at rest but insufficient during exercise, presenting clinically as intermittent claudication. Critical ischemia is produced when the reduction in blood flow results in a perfusion deficit at rest and is defined by the presence of pain at rest or trophic lesions in the legs. In this situation, precise diagnosis is fundamental, as there is a clear risk of loss of the limb if adequate blood flow is not re-established, either by surgery or by endovascular therapy. Differentiating between the two concepts is important in order to establish the therapeutic indication and the prognosis in patients with PAD.
Other causes include vasculitis and in situ thrombosis related to hypercoagulable states. Additional mechanisms of peripheral artery disease include arterial spasm and fibromuscular dysplasia. The cause and pathophysiology of arterial spasm are not fully understood, but it is hypothesised that they can occur secondary to trauma. The symptoms of claudication ensue when the artery spasms, or clamps down on itself, creating an obstruction. Similar to atherosclerosis, this leads to decreased blood flow to the tissue downstream of the obstruction. Thrombosis, or the formation of a blood clot, usually occurs due to stasis or trauma.
Diagnosis
Diagnosing or identifying peripheral artery disease requires a history of symptoms and a physical exam, followed by confirmatory testing. These tests could include CT scans (Computed Tomographic Angiography), MRA scans (Magnetic Resonance Angiography), or ultrasounds for imaging. In the setting of symptoms consistent with peripheral artery disease, a physician will then examine an individual for specific exam findings. Abnormal physical exam findings can lead a health care provider to consider a specific diagnosis. However, in order to confirm a diagnosis, confirmatory testing is required.
These findings are associated with peripheral artery disease:
Decreased or absent pulses
Muscle atrophy or wasting
Noticeable blueness of the affected limb
Decreased temperature (coolness) in the affected limb when compared to the other
Thickened nails
Smooth or shiny skin and hair loss
Buerger's test can check for pallor when the affected limb is in an elevated position. The limb is then moved from an elevated to a sitting position and checked for redness, which is called reactive hyperemia. Buerger's test is an assessment of arterial sufficiency, which is the ability of the artery to supply oxygenated blood to the tissue that it goes to.
Nonhealing lower extremity wound
If peripheral artery disease is suspected, the initial study is the ankle–brachial index (ABI). The ABI is a simple, non-invasive test that measures the ratio of systolic blood pressure in the ankle to the systolic blood pressure in the upper arm. This is based on the idea that if blood pressure readings in the ankle are lower than those in the arm, a blockage in the arteries that provide blood from the heart to the ankle is suspected. An ABI range of 0.90 to 1.40 is considered normal. A person is considered to have PAD when the ABI is ≤ 0.90. However, PAD can be further graded as mild to moderate if the ABI is between 0.41 and 0.90, and severe if the ABI is less than 0.40. These categories can provide insight into the disease course. Furthermore, ABI values of 0.91 to 0.99 are considered borderline, and values >1.40 indicate noncompressible arteries. If an ABI >1.40 is calculated, this could indicate vessel wall stiffness caused by calcification, which can occur in people with uncontrolled diabetes. Abnormally high ABIs (>1.40) are usually considered false negatives, and thus, such results merit further investigation and higher-level studies. Individuals with noncompressible arteries have an increased risk of cardiovascular mortality within a two-year period.
Individuals with suspected PAD with normal ABIs can undergo exercise testing for ABI. A baseline ABI is obtained prior to exercise. The patient is then asked to exercise (usually patients are made to walk on a treadmill at a constant speed) until claudication pain occurs (for a maximum of 5 minutes), after which the ankle pressure is again measured. A decrease in ABI of 15%–20% would be diagnostic of PAD.
If ABIs are abnormal, the next step is generally a lower limb Doppler ultrasound to look at the site of obstruction and extent of atherosclerosis. Other imaging can be performed by angiography, where a catheter is inserted into the common femoral artery and selectively guided to the artery in question. While injecting a radio-dense contrast agent, an X-ray is taken. Any blood flow-limiting blockage found in the X-ray can be identified and treated by procedures including atherectomy, angioplasty, or stenting. Contrast angiography is the most readily available and widely used imaging technique. Modern computerized tomography (CT) scanners provide direct imaging of the arterial system. Studies have shown the sensitivity and specificity of CT in identifying lesions with >50% stenosis to be 95% and 96%, respectively. As such, CT may be considered as an alternative to invasive angiography. An important distinction between the two is that, unlike invasive angiography, assessment of the arterial system with CT does not allow for vascular intervention.
Magnetic resonance angiography (MRA) is a noninvasive diagnostic procedure that uses a combination of a large magnet, radio frequencies, and a computer to produce detailed images of blood vessels inside the body. The advantages of MRA include its safety and ability to provide high-resolution, three-dimensional imaging of the entire abdomen, pelvis, and lower extremities in one sitting.
Classification
The two most commonly used methods to classify peripheral artery disease are the Fontaine and Rutherford systems of classification. The Fontaine stages were introduced by René Fontaine in 1954 to define the severity of chronic limb ischemia:
Stage I: asymptomatic
Stage IIa: intermittent claudication after walking a distance of more than 200 meters
Stage IIb: intermittent claudication after walking a distance of less than 200 meters
Stage III: rest pain
Stage IV: ulcers or gangrene of the limb
The Rutherford classification was created by the Society for Vascular Surgery and the International Society of Cardiovascular Surgery, introduced in 1986 and revised in 1997 (and known as the Rutherford classification after the lead author, Robert B. Rutherford). This classification system consists of four grades and seven categories (categories 0–6):
Grade 0, Category 0: asymptomatic
Grade I, Category 1: mild claudication
Grade I, Category 2: moderate claudication
Grade I, Category 3: severe claudication
Grade II, Category 4: rest pain
Grade III, Category 5: minor tissue loss; ischemic ulceration not exceeding ulcer of the digits of the foot
Grade IV, Category 6: major tissue loss; severe ischemic ulcers or frank gangrene
Moderate to severe PAD, classified by Fontaine's stages III to IV or Rutherford's categories 4 to 5, presents a limb threat (risk of limb loss) in the form of critical limb ischemia.
Recently, the Society for Vascular Surgery came out with a classification system based on "wound, ischemia and foot infection" (WIfI). This classification system, published in 2013, was created to account for the demographic changes that have occurred over the past forty years, including the increased incidence of high blood sugar and evolving techniques and abilities for revascularization. This system was created on the basis that ischemia and angiographic disease patterns are not the sole determinants of amputation risk. The WIfI classification system is broken up into two parts: wounds and ischemia. Wounds are graded 0 through 3 based on the presence of ulceration, gangrene, and ischemia.
Grade 0: no ulcer, no gangrene
Grade 1: small, shallow ulcer; no gangrene
Grade 2: deep ulcer with exposed tendon or bone, gangrene limited to toes
Grade 3: extensive, full-thickness ulcer; gangrene extending to the forefoot or midfoot
Ischemia is graded 0 through 3 based on ABI, ankle systolic pressure, and toe pressure.
Grade 0: ABI ≥0.80, ankle systolic pressure ≥100 mm Hg, toe pressure ≥60 mm Hg
Grade 1: arterial brachial index 0.6 to 0.79, ankle systolic pressure 70 to 100 mm Hg, toe pressure 40 to 59 mm Hg
Grade 2: ABI 0.4–0.59, ankle systolic pressure 50 to 70 mm Hg, toe pressure 30 to 39 mm Hg
Grade 3: ABI ≤0.39, ankle systolic pressure <50 mm Hg, toe pressure <30 mm Hg
The TASC (and TASC II) classification suggests PAD treatment is based on the severity of disease seen on an angiogram.
Screening
It is not clear if screening for disease in the general population is useful, as it has not been extensively studied. This includes screening with the ankle-brachial index (ABI), although a systematic review of the literature did not support the use of routine ABI screening in asymptomatic patients.
Testing for coronary artery disease or carotid artery disease is of unclear benefit. While PAD is a risk factor for abdominal aortic aneurysms (AAA), there is no data on screening individuals with asymptomatic PAD for abdominal aortic aneurysms. For people with symptomatic PAD, screening by ultrasound for AAA is not unreasonable.
Wearable devices and remote patient monitoring
A 2022 review found that a variety of wearable medical devices measuring different parameters (such as body temperature) were being combined with remote patient monitoring of PAD patients, in a goal to improve health outcomes.
Some studies propose the development of devices measuring oxygen continuously during exercise. This is because resting perfusion and metabolic activity are extremely low and differences between non-patients and PAD patients are barely measurable. As such, testing of vascular function and energetics requires a physiological challenge. Pulse oximeters can be inconvenient to wear during exercise and only give oxygen values at discrete time points, nor is there sufficient evidence to support any use in identifying PAD. Some publications and studies therefore discuss the use of wearable sensors measuring oxygen levels continuously in PAD patients, such as through transcutaneous means. However, because transcutaneous measurements are affected by movement (such as during exercise) and body temperature, use of oxygen sensors that are inserted subcutaneously as opposed to transcutaneously may most effectively help monitor a PAD patient’s progress and direct therapy decisions. To date, one oxygen sensing system has been approved for use in Europe to measure tissue perfusion in all PAD patients.
Treatment
Depending on the severity of the disease, these steps can be taken, according to these guidelines:
Lifestyle
Stopping smoking (cigarettes promote PAD and are a risk factor for cardiovascular disease)
Regular exercise for those with claudication helps open up alternative small vessels (collateral flow), and the limitation in walking often improves. Treadmill exercise (35 to 50 minutes, three or four times per week) has been reviewed as another treatment with a number of positive outcomes, including a reduction in cardiovascular events and improved quality of life. Supervised exercise programs increase pain-free walking time and the maximum walking distance in people with PAD.
Medication
Management of diabetes
Management of hypertension
Management of high cholesterol, and antiplatelet drugs such as aspirin and clopidogrel. Statins reduce clot formation and cholesterol levels, respectively, and can help with disease progression and address the other cardiovascular risks that the affected person is likely to have.
According to guidelines, taking aspirin or clopidogrel is recommended to reduce AMI ("heart attack"), stroke, and other causes of vascular death in people with symptomatic peripheral artery disease. It is recommended that aspirin and clopidogrel be taken alone and not in conjunction with one another (i.e., not as dual antiplatelet therapy). The recommended daily dosage of aspirin for treating PAD is between 75 and 325 mg, while the recommended daily dosage for clopidogrel is 75 mg. The effectiveness of both aspirin and clopidogrel to reduce the risk of cardiovascular ischemic events in people with symptomatic PAD is not well established. Research also suggests that low-dose rivaroxaban plus aspirin is effective as a new anti-thrombotic regimen for PAD.
Cilostazol can improve symptoms in some people. Pentoxifylline is of unclear benefit. Cilostazol may improve walking distance for people who experience claudication due to peripheral artery disease, but no strong evidence suggests that it improves the quality of life, decreases mortality, or decreases the risk of cardiovascular events.
Treatment with other drugs or vitamins is unsupported by clinical evidence, "but trials evaluating the effect of folate and vitamin B12 on hyperhomocysteinemia, a putative vascular risk factor, are near completion".
Revascularization
After a trial of the best medical treatment outlined above, if symptoms persist, patients may be referred to a vascular or endovascular surgeon. The benefit of revascularization is thought to correspond to the severity of ischemia and the presence of other risk factors for limb loss, such as wound and infection severity.
Angioplasty (or percutaneous transluminal angioplasty) can be done on solitary lesions in large arteries, such as the femoral artery, but may not have sustained benefits. Patency rates following angioplasty are highest for iliac arteries and decrease with arteries towards the toes. Other criteria that affect the outcome following revascularization are the length of the lesion and the number of lesions. There do not appear to be any long-term advantages or sustained benefits to placing a stent following angioplasty in order to hold the narrowing of the subsartorial artery open.
Atherectomy, in which the plaque is scraped off the inside of the vessel wall (albeit with no better results than angioplasty).
Vascular bypass grafting can be performed to circumvent a diseased area of the arterial vasculature. The great saphenous vein is used as a conduit if available, although artificial (Gore-Tex or PTFE) material is often used for long grafts when adequate venous conduit is unavailable.
When gangrene has set in, amputation may be required to prevent infected tissues from causing sepsis, a life-threatening illness.
Thrombolysis and thrombectomy are used in cases of arterial thrombosis or embolism.
shockwave intravascular lithotripsy, a minimally-invasive method which uses ultrasound waves to break up plaque within the artery without need for penetration. The method was first approved by the US Food and Drug Administration in February 2021, and has been used as a complement to more widely-used methods of atherectomy.
Guidelines
A guideline from the American College of Cardiology and American Heart Association for the diagnosis and treatment of lower extremity, renal, mesenteric, and abdominal aortic PAD was compiled in 2013, combining the 2005 and 2011 guidelines. For chronic limb-threatening ischemia, the ACCF/AHA guidelines recommend balloon angioplasty only for people with a life expectancy of 2 years or less or those who do not have an autogenous vein available. For those with a life expectancy greater than 2 years or who have an autogenous vein, bypass surgery is recommended.
Prognosis
Individuals with PAD have an "exceptionally elevated risk for cardiovascular events and the majority will eventually die of a cardiac or cerebrovascular etiology". Prognosis is correlated with the severity of the PAD as measured by an ABI. Large-vessel PAD increases mortality from cardiovascular disease significantly. PAD carries a greater than "20% risk of a coronary event in 10 years".
The risk is low that an individual with claudication will develop severe ischemia and require amputation, but the risk of death from coronary events is three to four times higher than matched controls without claudication. Of patients with intermittent claudication, only "7% will undergo lower-extremity bypass surgery, 4% major amputations, and 16% worsening claudication", but stroke and heart attack events are elevated, and the "5-year mortality rate is estimated to be 30% (versus 10% in controls)".
Epidemiology
The prevalence of PAD in the general population is 3–7%, affecting up to 20% of those over 70; 70%–80% of affected individuals are asymptomatic; only a minority ever require revascularization or amputation. Peripheral artery disease affects one in three diabetics over the age of 50. In the US, it affects 12–20 percent of Americans age 65 and older. Around 10 million Americans have PAD. Despite its prevalence and implications for cardiovascular risk, there are still low levels of awareness of risk factors and symptoms, with 26% of the population in the US reported to have knowledge of PAD.
In 2000, among people aged 40 years and older in the United States, rates of PAD were 4.3%. Rates were 14.5% for people aged 70 years or over. Within age groups, rates were generally higher for women than men. Non-Hispanic blacks had a rate of 7.9% compared to 4.4% in Non-Hispanic whites and 3.0% (1.4%–4.6%) in Mexican Americans.
The incidence of symptomatic PAD increases with age, from about 0.3% per year for men aged 40–55 years to about 1% per year for men aged over 75 years. The prevalence of PAD varies considerably depending on how PAD is defined and the age of the population being studied. People diagnosed with PAD have a greater risk of a MACE (Major Adverse Cardiac Event) and stroke. Their risk of developing a reinfarction, stroke, or transient ischemic attack within one year following a heart attack increases to 22.9%, compared to 11.4% for those without PAD.
The Diabetes Control and Complications Trial and the UK Prospective Diabetes Study trials in people with type 1 and type 2 diabetes, respectively, demonstrated that glycemic control is more strongly associated with microvascular disease than macrovascular disease. Pathologic changes occurring in small vessels may be more sensitive to chronically elevated glucose levels than atherosclerosis occurring in larger arteries.
Research
Research is being done on therapies to prevent the progression of PAD. In those who have developed critically poor blood flow to the legs, the benefit of autotransplantation of autologous mononuclear cells is unclear.
Only one randomized controlled trial has been conducted comparing vascular bypass to angioplasty for the treatment of severe PAD. The trial found no difference in amputation-free survival between vascular bypass and angioplasty at the planned clinical endpoint, but the trial has been criticized as being underpowered, limiting endovascular options, and comparing inappropriate endpoints. As of 2017, two randomized clinical trials are being conducted to better understand the optimal revascularization technique for severe PAD and critical limb ischemia (CLI), the BEST-CLI (Best Endovascular Versus Best Surgical Therapy for Patients With Critical Limb Ischemia) Trial and the BASIL-2 (Bypass Versus Angioplasty in Severe Ischaemia of the Leg – 2 )Trial.
In 2011, pCMV-vegf165 was registered in Russia as the first-in-class gene therapy drug for the treatment of PAD, including the advanced stage of critical limb ischemia.
| Biology and health sciences | Cardiovascular disease | Health |
489642 | https://en.wikipedia.org/wiki/Three-spined%20stickleback | Three-spined stickleback | The three-spined stickleback (Gasterosteus aculeatus) is a fish native to most inland and coastal waters north of 30°N. It has long been a subject of scientific study for many reasons. It shows great morphological variation throughout its range, ideal for questions about evolution and population genetics. Many populations are anadromous (they live in seawater but breed in fresh or brackish water) and very tolerant of changes in salinity, a subject of interest to physiologists. It displays elaborate breeding behavior (defending a territory, building a nest, taking care of the eggs and fry) and it can be social (living in shoals outside the breeding season) making it a popular subject of inquiry in fish ethology and behavioral ecology. Its antipredator adaptations, host-parasite interactions, sensory physiology, reproductive physiology, and endocrinology have also been much studied. Facilitating these studies is the fact that the three-spined stickleback is easy to find in nature and easy to keep in aquaria.
Evolution
The three-spined stickleback appears to be a rather old species that has remained morphologically unchanged for more than 10 million years. The oldest record of the species is from the Alta Mira Shale of the Monterey Formation of California, which preserves an articulated skeleton that appears essentially identical to the modern G. aculeatus complex. A slightly younger specimen is known from diatomite deposits from Lompoc that also belong to the Monterey Formation. Both specimens are known from marine deposits, suggesting a marine or anadramous lifestyle. The presence of the three-spined stickleback in the Miocene suggests that the three-spined stickleback complex must have diverged from the blackspotted stickleback prior to this point.
The iconic evolutionary traits of the three-spined stickleback, including rapid evolution in isolated environments, reduction of armor & pelvis, and ecological divisions into different niches, appear to have been a longstanding tendency of the Gasterosteus lineage. These traits are all present in a fossil stickleback, Gasterosteus doryssus, known from Miocene-aged freshwater deposits of the Truckee Formation in Nevada, US, which saw rapid evolution and morphological change over a period of several thousand years.
Description
This species can occasionally reach lengths of , but lengths of at maturity are more common. The body is laterally compressed. The base of the tail is slender. The caudal fin has 12 rays. The dorsal fin has 10–14 rays; in front of it are the three spines that give the fish its name (though some individuals may have only two or four). The third spine (the one closest to the dorsal fin) is much shorter than the other two. The back of each spine is joined to the body by a thin membrane. The anal fin has eight to 11 rays and is preceded by a short spine. The pelvic fins consist of just a spine and one ray. All spines can be locked in an erect position, making the fish extremely hard to swallow by a predator. The pectoral fins are large, with 10 rays. The body bears no scales, but is protected by bony plates on the back, flanks, and belly. Only one ventral plate is present, but the number of flank plates varies greatly across the distribution range and across habitat types (see below); it is normally higher in marine populations (some freshwater populations may in fact lack lateral plates altogether).
Dorsal coloration varies, but tends towards a drab olive or a silvery green, sometimes with brown mottling. The flanks and belly are silvery. In males during the breeding season, the eyes become blue and the lower head, throat, and anterior belly turn bright red. The throat and belly of breeding females can turn slightly pink. A few populations, however, have breeding males which are all black or all white.
Habitat and distribution
The three-spined stickleback is found only in the Northern Hemisphere, where it usually inhabits coastal waters or freshwater bodies. It can live in either fresh, brackish, or salt water. It prefers slow-flowing water with areas of emerging vegetation. It can be found in ditches, ponds, lakes, backwaters, quiet rivers, sheltered bays, marshes, and harbours.
In North America, it ranges along the East Coast from Chesapeake Bay to the southern half of Baffin Island and the western shore of Hudson Bay, and along the West Coast from southern California to the western shore of Alaska and the Aleutian Islands. It can be found throughout Europe between 35 and 70°N. In Asia, the distribution stretches from Japan and the Korean peninsula to the Bering Straits.
Its distribution could be said to be circumpolar were it not for the fact that it is absent from the north coast of Siberia, the north coast of Alaska, and the Arctic islands of Canada.
Variation in morphology and distribution
Three subspecies are currently recognized by the IUCN:
G. a. aculeatus is found in most of the species range, and is the subspecies most strictly termed the three-spined stickleback; its common name in Britain is the tiddler, although "tittlebat" is also sometimes used.
G. a. williamsoni, the unarmored threespine stickleback, is found only in North America; its recognised range is southern California, though isolated reports have been made of it occurring in British Columbia and Mexico;
G. a. santaeannae, the Santa Ana stickleback, is also restricted to North America.
These subspecies actually represent three examples from the enormous range of morphological variation present within three-spined sticklebacks. Hybrids between some of these morphs show foraging disadvantages, a form of reinforcement in the course of speciation. This is evidence for speciation by reinforcement.
Overall these morphs fall into two rough categories, the anadromous and the freshwater forms:
The anadromous form spends most of its adult life eating plankton and fish in the sea, and returns to freshwater to breed. The adult fish are typically between 6 and 10 cm long, and have 30 to 40 lateral armour plates along their sides. They also have long dorsal and pelvic spines. The anadromous form is morphologically similar all around the Northern Hemisphere, such that anadromous fish from the Baltic, the Atlantic and the Pacific all resemble each other quite closely.
Three-spined stickleback populations are also found in freshwater lakes and streams. These populations were probably formed when anadromous fish started spending their entire lifecycle in fresh water, and thus evolved to live there all year round. Freshwater populations are extremely morphologically diverse, to the extent that many observers (and some taxonomists) would describe a new subspecies of three-spined stickleback in almost every lake in the Northern Hemisphere. One consistent difference between freshwater populations and their anadromous ancestors is the amount of body armour, as the majority of freshwater fish only have between none and 12 lateral armor plates, and shorter dorsal and pelvic spines. However, also large morphological differences occur between lakes. One major axis of variation is between populations found in deep, steep-sided lakes and those in small, shallow lakes. The fish in the deep lakes typically feed in the surface waters on plankton, and often have large eyes, with short, slim bodies and upturned jaws. Some researchers refer to this as the limnetic form. Fish from shallow lakes feed mainly on the lake bed, and are often long and heavy bodied with relatively horizontal jaws and small eyes. These populations are referred to as the benthic form.
Since each watershed was probably colonised separately by anadromous sticklebacks, morphologically similar populations in different watersheds or on different continents are widely believed to have evolved independently. A unique population is found in the meromictic Pink Lake in Gatineau Park, Quebec. Populations have been observed rapidly adapting to different conditions, such as in Lake Union, where sticklebacks have lost and regained armor plates in response to pollution from human activity around the watershed.
One aspect of this morphological variation is that a number of lakes contain both a limnetic and a benthic type, and these do not interbreed with each other. Evolutionary biologists often define species as populations that do not interbreed with each other (the biological species concept), thus the benthics and limnetics within each lake would constitute separate species. These species pairs are an excellent example of how adaptation to different environments (in this case feeding in the surface waters or on the lake bed) can generate new species. This process has come to be termed ecological speciation. This type of species pair is found in British Columbia. The lakes themselves only contain three-spined sticklebacks and cutthroat trout, and all are on islands. Tragically, the pair in Hadley Lake on Lasqueti Island was destroyed in the mid-1980s by the introduction of a predatory catfish, and the pair in Enos Lake on Vancouver Island has started to interbreed and are no longer two distinct species. The two remaining pairs are on Texada Island, in Paxton Lake and Priest Lake, and they are listed as Endangered in the Canadian Species at Risk Act.
Other species pairs which consist of a well-armored marine form and a smaller, unarmored freshwater form are being studied in ponds and lakes in south-central Alaska that were once marine habitats such as those uplifted during the 1964 Alaska earthquake. The evolutionary dynamics of these species pairs are providing a model for the processes of speciation which has taken place in less than 20 years in at least one lake. In 1982, a chemical eradication program intended to make room for trout and salmon at Loberg Lake, Alaska, killed the resident freshwater populations of sticklebacks. Oceanic sticklebacks introduced through nearby Cook Inlet recolonized the lake. In just 12 years beginning in 1990, the frequency of the oceanic form dropped steadily, from 100% to 11%, while a variety with fewer plates increased to 75% of the population, with various intermediate forms making up another small fraction. This rapid evolution is thought to be possible through genetic variations that confer competitive advantages for survival in fresh water when conditions shift rapidly from salt to fresh water. However, the actual molecular basis of this evolution still remains unknown.
Although sticklebacks are found in many locations around the coasts of the Northern Hemisphere and are thus viewed by the IUCN as species of least concern, the unique evolutionary history encapsulated in many freshwater populations indicates further legal protection may be warranted.
Diet
In its different forms or stages of life, the three-spined stickleback can be a bottom-feeder (most commonly chironomid larvae and amphipods) or a planktonic feeder in lakes or in the ocean; it can also consume terrestrial prey fallen to the surface. It can cannibalize eggs and fry.
Life history
Many populations take two years to mature and experience only one breeding season before dying, and some can take up to three years to reach maturity. However, some freshwater populations and populations at extreme latitudes can reach maturity in only one year.
Reproduction
Sexual maturation depends on environmental temperature and photo-period. Longer days and warmer days stimulate brighter colouration in males and the development of eggs in females.
From late April, males and females move from deeper waters to shallow areas. There, each male defends a territory where he builds a nest on the bottom. He starts by digging a small pit. He then fills it with plant material (often filamentous algae), sand, and various debris which he glues together with spiggin, a proteinaceous substance secreted from the kidneys. The word spiggin is derived from spigg, the Swedish name for the three-spined stickleback. He then creates a tunnel through the more or less spherical nest by swimming vigorously through it. Nest building typically takes 5–6 hours though it may also be spread out over several days. After this, the male courts gravid females that pass by with a zigzag dance. (In some populations, the male leads the female to the nest, rather than doing the zigzag dance.) He approaches a female by swimming very short distances left and right, and then swims back to the nest in the same way. If the female follows, the male often pokes his head inside the nest, and may swim through the tunnel. The female then swims through the tunnel as well, where she deposits 40–300 eggs. The male follows to fertilize the eggs. The female is then chased away by the male. For the duration of the eggs' development, the male will chase away other males and non-gravid females. He may, however, court other gravid females (more than one batch of eggs can be deposited in the same nest).
The sequence of territorial courtship and mating behaviours was described in detail by Niko Tinbergen in a landmark early study in ethology. Tinbergen showed that the red colour on the throat of the territorial male acts as a simple sign stimulus, releasing aggression in other males and attracting females. The red colouration may also be used by females as a way to assess male quality. Red colouration is produced from carotenoids found in the diet of the fish. As carotenoids cannot be synthesised de novo, the degree of colouration gives an indication of male quality (ability to find food), with higher-quality males showing more intense colouration. Also, males that bear fewer parasites tend to exhibit brighter red colours. Many studies have shown that females prefer males with brighter red colouration. However, the response to red is not universal across the entire species, with black throated populations often found in peat-stained waters.
The male takes care of the developing eggs by fanning them. He lines himself up with the entrance of the nest tunnel and swims on the spot. The movement of his pectoral fins creates a current of water through the nest, bringing fresh (well-oxygenated) water to the eggs. He does this not only during the day, but throughout the night, as well. Fanning levels tend to increase until the eggs are about to hatch, which takes 7–8 days at 18–20 °C. Fanning levels also increase when the water is poorly oxygenated. Towards the end of the egg development phase, the male often makes holes in the roof and near the rim of the nest, presumably to improve ventilation of the nest during fanning at a time when the eggs are more metabolically active. Once the young hatch, the male attempts to keep them together for a few days, sucking up any wanderers into his mouth and spitting them back into the nest. Afterwards, the young disperse and the nest is either abandoned by the male, or repaired in preparation for another breeding cycle.
In Nova Scotia, a form of three-spined stickleback departs from the usual pattern of parental care. Unlike other sticklebacks that nest on the substrate, Nova Scotian male sticklebacks build nests in mats of filamentous algae. Surprisingly, almost immediately after fertilization, the males disperse the eggs from the nest and resume soliciting females for eggs. Hence, there appears to have been a loss of parental care in this population. Because these males have reduced dorsal pigmentation, resulting a pearlescent white appearance, they have been dubbed "white sticklebacks". It is currently unknown whether they are a distinct species, or simply a morph of the common Atlantic stickleback.
As the breeding cycle of the three-spined stickleback is light and temperature dependent, it is also possible to manipulate breeding in the lab. For example, it is possible to stimulate sticklebacks to breed twice in a calendar year, instead of once, under the right conditions. This can be useful for genetic and behavioural multi-generational studies.
Infection with the cestode parasite Schistocephalus solidus can cause a reduction in egg mass or complete absence of eggs in female three-spined sticklebacks.
Cooperative behavior
Some evidence indicates the existence of cooperative behavior among three-spined sticklebacks, mainly cooperative predator inspection. Predator inspection appears to allow acquisition of information about the risk a potential predator presents, and may deter attack, with the cost being an increased chance of being attacked if the predator proves to be hungry.
Tit for tat strategy
Sticklebacks are known to cooperate in a tit for tat (TFT) strategy when doing predator inspection. The idea behind TFT is that an individual cooperates on the first move and then does whatever its opponent does on the previous move. This allows for a combination of collaborative (it starts by cooperating), retaliatory (punishes defection), and forgiving (respond to cooperation of others, even if they had defected previously) behavioral responses. When three-spined sticklebacks approaching a live predator were provided with either a simulated cooperating companion or a simulated defecting one, the fish behaved according to tit-for-tat strategy, supporting the hypothesis that cooperation can evolve among egoists.
Typically, sticklebacks operate in pairs. Individuals have partners with which they repeatedly perform pairwise predator inspection visits. Two reciprocal pairs per trial occur significantly more often than what was expected due to chance. These results provide further evidence for a tit-for-tat cooperation strategy in sticklebacks.
Stickleback behavior is often cited as an archetypal example of cooperative behavior during predator inspection. Fish from three sites differing in predation risk inspected a model predator in pairs and reciprocated both cooperative moves and defections by the partner, but not on every opportunity. Sticklebacks that originated in the two sites containing piscivorous fish were more likely to reciprocate following a cooperative move than following a defection. Individuals from higher-risk sites were generally more cooperative. Individuals accompanied by a model companion show reciprocal moves of cooperation and defection in response to the model's movements about a third of the time. Both examples of stickleback behavior demonstrate the elements of a strategy of cooperation that may resemble tit-for-tat.
Partner-dependence
The tit-for-tat cooperation strategy has been shown to be evident in sticklebacks. In addition, the size of a stickleback's partner fish may also be a factor in determining what a stickleback will do when both fish are faced with a predator. Two sticklebacks simultaneously presented to a rainbow trout, a predator much larger in size, will have differing risks of being attacked. Usually, the larger of the two sticklebacks has a higher risk of being attacked. Individual sticklebacks are more likely to move closer to a trout (or some other predator) when a larger potential partner moves close to the trout than when a smaller partner approaches the trout. Although both large and small partners behave similarly, a small partner's behavior affects the strategy of the test fish more than that of the large partner. Regardless of whether it is alone or with a partner that cooperates, a larger fish will approach a predator more closely than does a smaller fish. If a partner defects, then a stickleback's condition-factor (i.e. its ability to flee) determines how closely it approaches the predator rather than the stickleback's size. Both the strategy and reaction to different-sized partners seem to be dependent on whether the partner cooperates or defects.
Parasites
The three-spined stickleback is a secondary intermediate host for the hermaphroditic parasite Schistocephalus solidus, a tapeworm of fish and fish-eating birds. The tapeworm passes into sticklebacks through its first intermediate hosts, cyclopoid copepods, when these are eaten by the fish. The parasite matures into its third larval stage, the plerocercoid, in the abdomen of the stickleback. Infected sticklebacks are afterwards consumed by fish-eating birds, which serve as the tapeworm's definitive host.
Another common parasite of the three-spined stickleback is the microsporidian Glugea anomala. Naturally infections with G. anomala lose weight compared to uninfected individuals, but do not cause size differences between individuals. Glugea anomala also correlates behavioural changes, such as increased shoaling, increased sociability and activity, and reduced boldness. It is unknown whether these differences in behaviour are due to certain personality traits predisposing individuals to infections, or whether infections change behaviour.
Genetics
Three-spined sticklebacks have recently become a major research organism for evolutionary biologists trying to understand the genetic changes involved in adapting to new environments. The entire genome of a female fish from Bear Paw Lake in Alaska was recently sequenced by the Broad Institute and many other genetic resources are available. This population is under risk from the presence of introduced northern pike in a nearby lake. Three-spined sticklebacks are also used for researching sex-specific brain gene expression. Parents exposed to predator models produced offspring with different gene expressions compared to those that were not exposed to predators. Non-overlapping genes appear highly influenced by the sex of the parent, with genes being differentially expressed in offspring based on whether the male or female parent was exposed to predation.
Eco-evolutionary dynamics
Three-spined stickleback research has been central to the field of eco-evolutionary dynamics. Eco-evolutionary dynamics is an area of study investigating how ecological processes (e.g., population dynamics, community interactions, and nutrient cycling) affect how populations evolve, and in turn, how these patterns of evolution feed back to affect ecological processes. Importantly, these dynamics arise when substantial evolutionary change occurs on the same time scale as ecological change (i.e., less than 1,000 generations). Three-spined stickleback are particularly useful for studying eco-evolutionary dynamics because multiple populations have evolved rapidly and in predictable, repeated patterns after colonizing new environments. These repeated patterns of evolution allow scientists to assess whether the impacts of stickleback evolution on ecological processes are reproducible.
An eco-evolutionary framework has been used to explore multiple aspects of stickleback biology. Notably, this research has focused on how populations of three-spined stickleback have diverged to occupy different ecological niches (a process called adaptive radiation) and how sticklebacks have coevolved with their parasites.
Eco-evolutionary dynamics of adaptive radiation
Most eco-evolutionary dynamics research in sticklebacks has focused on how the adaptive radiation of different ecotypes affects ecological processes. Ecotypes represent genetically and morphologically recognizable populations occupying distinct ecological niches. In three-spined stickleback, divergent ecotypes are often found as sympatric (i.e., co-occurring) or parapatric (i.e., partially overlapping, but mostly isolated) species pairs, including benthic—limnetic pairs, freshwater—anadromous pairs, and lake—stream pairs. Pairs of stickleback ecotypes have diverged at time scales ranging from 10,000 years to only decades ago.
Different combinations of stickleback ecotypes affect ecosystem processes in different ways. For example, the combined presence of specialized benthic and limnetic sticklebacks has a different effect on the diversity and abundance of prey species compared to the presence of only a generalist ancestral stickleback ecotype. Notably, this effect appears to be driven by limnetic sticklebacks specializing on zooplankton prey, rather than by an increase in the number of co-occurring stickleback species alone. The impacts of ecotype specialization on prey communities can even affect the abundance of algae and cyanobacteria that do not directly interact with sticklebacks, along with aspects of the abiotic environment, such as the amount of ambient light available for photosynthesis and levels of dissolved oxygen, carbon, and phosphorus. These diverse changes in ecosystem processes can persist to affect natural selection on subsequent stickleback generations, potentially shaping how stickleback populations will evolve in the future. Because the presence of specialist verses generalist ecotypes can impact ecosystems in a way that, in turn, affects selection on future stickleback generations, the adaptive radiation of specialized ecotypes could drive eco-evolutionary feedback loops in natural populations.
Eco-evolutionary dynamics of host-parasite interactions
Sticklebacks have also been studied to investigate the eco-evolutionary dynamics of host-parasite coevolution. Three-spined sticklebacks can be hosts to a variety of parasites (e.g., Schistocephalus solidus, a common tapeworm of fish and fish-eating birds). The diversity of parasite species within individual stickleback is influenced by an individual's dietary niche and immune response. This covariation between parasite infection and host traits is likely a consequence of eco-evolutionary feedback, whereby the evolution of dietary and parasite resistance traits in sticklebacks alters parasite reproduction and infection rates, which in turn affects parasite exposure and selection on parasite resistance in sticklebacks. These feedbacks can also extend beyond stickleback-parasite interactions to modify ecosystem processes. Specifically, differences in resistance and infection rates among stickleback ecotypes can alter how sticklebacks affect the abundance of prey species and levels of dissolved nutrients and oxygen. These ecosystem impacts can further affect selection on sticklebacks in subsequent generations, which suggests a complex feedback loop between the evolution of host-parasite interactions, community composition, and abiotic conditions.
Common methods
Many researchers have used mesocosm experiments to test how the adaptive radiation of stickleback ecotypes and stickleback-parasite interactions can impact ecological processes. In these experiments, researchers simulate the natural environments of sticklebacks in enclosed tanks, including natural plant and invertebrate communities and freshwater ecological zones. They then systematically manipulate an independent variable (e.g., which stickleback ecotypes were present or the presence of parasites), and measured differences in biotic and abiotic aspects of ecosystems among the different stickleback treatments. In some cases, researchers have then tested for potential feedback loops between ecotype evolution and ecological change by removing the adult stickleback from the mesocosms and replacing them with juveniles of different ecotypes. By doing so, the researchers could then measure how the effects of adult sticklebacks on their ecosystems influenced overall juvenile fitness (e.g., survival and growth rates) and differences in fitness between juveniles of different ecotypes.
| Biology and health sciences | Acanthomorpha | Animals |
489658 | https://en.wikipedia.org/wiki/Nd%3AYAG%20laser | Nd:YAG laser | Nd:YAG (neodymium-doped yttrium aluminum garnet; Nd:Y3Al5O12) is a crystal that is used as a lasing medium for solid-state lasers. The dopant, neodymium in the +3 oxidation state, Nd(III), typically replaces a small fraction (1%) of the yttrium ions in the host crystal structure of the yttrium aluminum garnet (YAG), since the two ions are of similar size. It is the neodymium ion which provides the lasing activity in the crystal, in the same fashion as red chromium ion in ruby lasers.
Laser operation of Nd:YAG was first demonstrated by J.E. Geusic et al. at Bell Laboratories in 1964.
Technology
Nd:YAG lasers are optically pumped using a flashtube or laser diodes. These are one of the most common types of laser, and are used for many different applications.
Nd:YAG lasers typically emit light with a wavelength of 1064 nm, in the infrared. However, there are also transitions near 946, 1120, 1320, and 1440 nm. Nd:YAG lasers operate in both pulsed and continuous mode. Pulsed Nd:YAG lasers are typically operated in the so-called Q-switching mode: An optical switch is inserted in the laser cavity waiting for a maximum population inversion in the neodymium ions before it opens. Then the light wave can run through the cavity, depopulating the excited laser medium at maximum population inversion. In this Q-switched mode, output powers of 250 megawatts and pulse durations of 10 to 25 nanoseconds have been achieved. The high-intensity pulses may be efficiently frequency doubled to generate laser light at 532 nm, or higher harmonics at 355, 266 and 213 nm.
Nd:YAG absorbs mostly in the bands between 730–760 nm and 790–820 nm. At low current densities krypton flashlamps have higher output in those bands than do the more common xenon lamps, which produce more light at around 900 nm. The former are therefore more efficient for pumping Nd:YAG lasers.
The amount of the neodymium dopant in the material varies according to its use. For continuous wave output, the doping is significantly lower than for pulsed lasers. The lightly doped CW rods can be optically distinguished by being less colored, almost white, while higher-doped rods are pink-purplish.
Other common host materials for neodymium are: YLF (yttrium lithium fluoride, 1047 and 1053 nm), YVO4 (yttrium orthovanadate, 1064 nm), and glass. A particular host material is chosen in order to obtain a desired combination of optical, mechanical, and thermal properties. Nd:YAG lasers and variants are pumped either by flashtubes, continuous gas discharge lamps, or near-infrared laser diodes (DPSS lasers). Prestabilized laser (PSL) types of Nd:YAG lasers have proved to be particularly useful in providing the main beams for gravitational wave interferometers such as LIGO, VIRGO, GEO600 and TAMA.
Applications
Medicine
Nd:YAG lasers are used in ophthalmology to correct posterior capsular opacification, after cataract surgery, for peripheral iridotomy in patients with chronic and acute angle-closure glaucoma, where it has largely superseded surgical iridectomy, for the treatment of vitreous eye floaters, for pan-retinal photocoagulation in the treatment of proliferative diabetic retinopathy, and to damage the retina in ophthalmology animal research.
Nd:YAG lasers emitting light at 1064 nm have been the most widely used laser for laser-induced thermotherapy, in which benign or malignant lesions in various organs are ablated by the beam.
In oncology, Nd:YAG lasers can be used to remove skin cancers. They are also used to reduce benign thyroid nodules, and to destroy primary and secondary malignant liver lesions.
To treat benign prostatic hyperplasia (BPH), Nd:YAG lasers can be used for laser prostate surgery—a form of transurethral resection of the prostate.
These lasers are also used extensively in the field of cosmetic medicine for laser hair removal and the treatment of minor vascular defects such as spider veins on the face and legs. Nd:YAG lasers are also used to treat venous lake lip lesions. Recently Nd:YAG lasers have been used for treating dissecting cellulitis of the scalp, a rare skin disease.
Using hysteroscopy the Nd:YAG laser has been used for removal of uterine septa within the inside of the uterus.
In podiatry, the Nd:YAG laser is being used to treat onychomycosis, which is fungus infection of the toenail. The merits of laser treatment of these infections are not yet clear, and research is being done to establish effectiveness.
Dentistry
Nd:YAG dental lasers have been used for the removal of dental caries as an alternative to drill therapy, although evidence supporting its use is of low quality. They have also been used for soft tissue surgeries in the oral cavity, such as gingivectomy, periodontal sulcular debridement, LANAP, and pulpotomy. Nd:YAG dental lasers have also been shown to be effective at treating and preventing dental hypersensitivity, as an adjunct for periodontal instrumentation, and for the treatment of recurrent aphthous stomatitis.
Manufacturing
Nd:YAG lasers are used in manufacturing for engraving, etching, or marking a variety of metals and plastics, or for metal surface enhancement processes like laser peening. They are extensively used in manufacturing for cutting and welding steel, semiconductors and various alloys. For automotive applications (cutting and welding steel) the power levels are typically 1–5 kW. Super alloy drilling (for gas turbine parts) typically uses pulsed Nd:YAG lasers (millisecond pulses, not Q-switched). Nd:YAG lasers are also employed to make subsurface markings in transparent materials such as glass or acrylic glass and in white and transparent polycarbonate for identity documents. Lasers of up to 2 kW are used for selective laser melting of metals in additive layered manufacturing. In aerospace applications, they can be used to drill cooling holes for enhanced air flow/heat exhaust efficiency.
Nd:YAG lasers are also used in the non-conventional rapid prototyping process laser engineered net shaping (LENS).
Laser peening typically uses a high energy (10 to 40 joule) 10 to 30 nanosecond pulse. The laser beam is focused down to a few millimeters in diameter to deposit gigawatts of power on the surface of a part. Laser peening is unlike other manufacturing processes in that it neither heats nor adds material; it is a mechanical process of cold working the metallic component to impart compressive residual stresses. Laser peening is widely used in gas-fired turbine engines in both aerospace and power generation to increase strength and improve resistance to damage and metal fatigue.
Fluid dynamics
Nd:YAG lasers can be used for flow visualization techniques in fluid dynamics (for example particle image velocimetry or laser-induced fluorescence).
Biophysics
Nd:YAG lasers are frequently used to build optical tweezers for biological applications. This is because Nd:YAG lasers mostly emit at a wavelength of 1064 nm. Biological samples have a low absorption coefficient at this wavelength, as biological samples are usually mostly made up of water.
As such, using an Nd:YAG laser minimizes the damage to the biological sample being studied.
Automotive
Researchers from Japan's National Institutes of Natural Sciences are developing laser igniters that use YAG chips to ignite fuel in an engine, in place of a spark plug. The lasers use several 800 picosecond long pulses to ignite the fuel, producing faster and more uniform ignition. The researchers say that such igniters could yield better performance and fuel economy, with fewer harmful emissions.
Military
The Nd:YAG laser is the most common laser used in laser designators and laser rangefinders.
During the Iran–Iraq War, Iranian soldiers suffered more than 4000 cases of laser eye injury, caused by a variety of Iraqi sources including tank rangefinders. The 1064 nm wavelength of Nd:YAG is thought to be particularly dangerous, as it is invisible and initial exposure is painless.
The Chinese ZM-87 blinding laser weapon uses a laser of this type, though only 22 have been produced due to their prohibition by the Convention on Certain Conventional Weapons. North Korea is reported to have used one of these weapons against American helicopters in 2003.
Cavity ring-down spectroscopy (CRDS)
The Nd:YAG may be used in the application of cavity ring-down spectroscopy, which is used to measure the concentration of some light-absorbing substance.
Laser-induced breakdown spectroscopy (LIBS)
A range of Nd:YAG lasers are used in analysis of elements in the periodic table. Though the application by itself is fairly new with respect to conventional methods such as XRF or ICP, it has proven to be less time consuming and a cheaper option to test element concentrations. A high-power Nd:YAG laser is focused onto the sample surface to produce plasma. Light from the plasma is captured by spectrometers and the characteristic spectra of each element can be identified, allowing concentrations of elements in the sample to be measured.
Laser pumping
Nd:YAG lasers, mainly via their second and third harmonics, are widely used to excite dye lasers either in the liquid or solid state. They are also used as pump sources for vibronically broadened solid-state lasers such as Cr4+:YAG or via the second harmonic for pumping Ti:sapphire lasers.
Additional frequencies
For many applications, the infrared light is frequency-doubled or -tripled using nonlinear optical materials such as lithium triborate to obtain visible (532 nm, green) or ultraviolet light. Cesium lithium borate generates the 4th and 5th harmonics of the Nd:YAG 1064 nm fundamental wavelength. A green laser pointer is a frequency doubled Nd:YVO4 diode-pumped solid state laser (DPSS laser). Nd:YAG can be also made to lase at its non-principal wavelength. The line at 946 nm is typically employed in "blue laser pointer" DPSS lasers, where it is doubled to 473 nm.
Physical and chemical properties of Nd:YAG
Properties of YAG crystal
Formula: Y3Al5O12
Molecular weight: 596.7
Crystal structure: Cubic
Hardness: 8–8.5 (Mohs)
Melting point: 1970 °C (3540 °F)
Density: 4.55 g/cm3
Refractive index of Nd:YAG
Properties of Nd:YAG @ 25 °C (with 1% Nd doping)
Formula: Y2.97Nd0.03Al5O12
Weight of Nd: 0.725%
Atoms of Nd per unit volume: 1.38×1020 /cm3
Charge state of Nd: 3+
Emission wavelength: 1064 nm
Transition: 4F3/2 → 4I11/2
Duration of fluorescence: 230 μs
Thermal conductivity: 0.14 W·cm−1·K−1
Specific heat capacity: 0.59 J·g−1·K−1
Thermal expansion: 6.9×10−6 K−1
dn/dT: 7.3×10−6 K−1
Young's modulus: 3.17×104 K·g/mm−2
Poisson's ratio: 0.25
Resistance to thermal shock: 790 W·m−1
| Technology | Lasers | null |
490225 | https://en.wikipedia.org/wiki/Chimaera | Chimaera | Chimaeras are cartilaginous fish in the order Chimaeriformes (), known informally as ghost sharks, rat fish, spookfish, or rabbit fish; the last three names are not to be confused with rattails, Opisthoproctidae, or Siganidae, respectively.
At one time a "diverse and abundant" group (based on the fossil record), their closest living relatives are sharks and rays, though their last common ancestor with them lived nearly 400 million years ago. Living species (aside from plough-nose chimaeras) are largely confined to deep water.
Anatomy
Chimaeras are soft-bodied, shark-like fish with bulky heads and long, tapered tails; measured from the tail, they can grow up to in length. Like other members of the class Chondrichthyes, chimaera skeletons are entirely cartilaginous, or composed of cartilage. Males use forehead denticles to grasp a female by a fin during copulation. The gill arches are condensed into a pouch-like bundle covered by a sheet of skin (an operculum), with a single gill-opening in front of the pectoral fins.
The pectoral fins are large enough to generate lift at a relaxed forward momentum, giving the chimaera the appearance of "flying" through the water. Further back on the body are also a pair of smaller pelvic fins, and some genera bear an anal fin in front of the tail. In chimaerids and rhinochimaerids, the tail is leptocercal, meaning that it is thin and whip-like, edged from above and below by fins of similar size. In callorhinchids, the tail is instead heterocercal, with a larger upper lobe inclined upwards, similar to many sharks. There are two dorsal fins: a large triangular first dorsal fin and a low rectangular or depressed second dorsal fin. For defense, some chimaeras have a venomous spine on the front edge of the dorsal fin.
In many species, the bulbous snout is modified into an elongated sensory organ, capable of electroreception to find prey. The cartilaginous skull is holostylic, meaning that the palatoquadrate (upper jaw cartilage) is completely fused to the neurocranium (cranial cartilage). This contrasts with modern sharks, where the palatoquadrate is movable and detachable, a trait known as hyostyly. The back of the head is supported by a complex of fused vertebrae called the synarcual, which also connects to the dorsal fin spine.
Instead of sharks' many sharp, consistently-replaced teeth, chimaeras have just six large, permanent tooth-plates, which grow continuously throughout their entire life. These tooth-plates are arranged in three pairs, with one pair at the tip of the lower jaws and two pairs along the upper jaws. They together form a protruding, beak-like crushing and grinding mechanism, comparable to the incisor teeth of rodents and lagomorphs (hence the name "rabbit fish"). Chimaera teeth are unique among vertebrates, due to their mode of mineralization. Most of each plate is formed by relatively soft osteodentin, but the active edges are supplemented by a unique hypermineralized tissue called pleromin. Pleromin is an extremely hard enamel-like tissue, arranged into sheets or beaded rods, but it is deposited by mesenchyme-derived cells similar to those that form bone. In addition, pleromin's hardness is due to the mineral whitlockite, which crystalizes within the teeth as the animal matures. Other vertebrates with hypermineralized teeth rely on enamel, which is derived from ameloblasts and encases round crystals of the mineral apatite.
Chimaeras also differ from sharks in that they have separate anal and urogenital openings.
Behavior
Chimaeras live in temperate ocean floors, with some species inhabiting depths exceeding , with relatively few modern species regularly inhabiting shallow water. Exceptions include the members of the genus Callorhinchus, the rabbit fish and the spotted ratfish, which locally or periodically can be found at shallower depths. Consequently, these are also among the few species kept in public aquaria. They live in all the oceans except for the Arctic and Antarctic oceans.
Diet
The usual diet of chimaeras consist of crustaceans, and more specifically, they include ophiurans and molluscs. Modern species are demersal durophages, but they used to be more diverse. The Carboniferous period had forms that lived as specialised suction feeders in the water column.
Reproduction
Chimaera reproduction resembles that of sharks in some ways: males employ claspers for internal fertilization of females and females lay eggs within spindle-shaped, leathery egg cases.
Unlike sharks, male chimaeras have retractable sexual appendages (known as tentacula) to assist mating. The frontal tentaculum, a bulbous rod which extends out of the forehead, is used to clutch the females' pectoral fins during mating. The prepelvic tentacula are serrated hooked plates normally hidden in pouches in front of the pelvic fins, and they anchor the male to the female. Lastly, the pelvic claspers (sexual organs shared by sharks) are fused together by a cartilaginous sheathe before splitting into a pair of flattened lobes at their tip.
Parasites
As other fish, chimaeras have a number of parasites. Chimaericola leptogaster (Chimaericolidae) is a monogenean parasite of the gills of Chimaera monstrosa; the species can attain in length.
Conservation and threats
Despite their secluded habits, some chimaera species may be threatened by overfishing through bycatch or commercial exploitation. No species are listed as Endangered according to the IUCN, but four are listed as Vulnerable, four more as Near Threatened, and many more as Data Deficient (too rare to evaluate). Many species have restricted ranges and practically none have had their movement patterns studied. In addition, bycatch reports are usually insufficiently precise to the species or even genus level, so it is difficult to keep track of bycatch on a species-by-species basis. This lack of data renders chimaera species especially susceptible to overlooked population declines.
Several near-shore species are purposefully caught for their meat, especially callorhinchids, Hydrolagus bemisi (pale ghost shark), and Hydrolagus novaezealandiae (dark ghost shark). Modern quotas have helped to moderate collection of these species to a sustainable level, though Callorhinchus milii (the Australian ghostshark) experienced severe overfishing in the 20th century before protections were enacted. Neoharriotta pinnata (sicklefin chimaera) is targeted along the coast of India for its liver oil, and a recent decline of catch rates may indicate a population crash. Even species without commercial exploitation can fall victim to bycatch: Callorhinchus callorynchus (American elephantfish), Neoharriotta carri (dwarf sicklefin chimaera), Chimaera monstrosa (rabbit fish), Chimaera ogilbyi (Ogilby's ghostshark), Hydrolagus colliei (spotted ratfish), and Hydrolagus melanophasma (eastern Pacific black ghostshark) all have bycatch rates exceeding 10% in certain parts of their range, and some are experiencing steep declines. Chimaeras have mostly avoided harvesting for the fin trade, which threatens many true sharks.
Another threat is habitat destruction of coastal nurseries (by urban development) or deepwater reefs (by deep sea mining and trawling). Near-shore species such as Callorhinchus milii are vulnerable to the effects of climate change: stronger storms and warmer seawater are predicted to increase egg mortality by disrupting the stable environments necessary to complete incubation.
Classification
In some classifications, the chimaeras are included (as subclass Holocephali) in the class Chondrichthyes of cartilaginous fishes; in other systems, this distinction may be raised to the level of class. Chimaeras also have some characteristics of bony fishes.
A renewed effort to explore deep water and to undertake taxonomic analysis of specimens in museum collections led to a boom during the first decade of the 21st century in the number of new species identified. A preliminary study found 8% of species to be threatened. There are over 50 extant species in six genera and three families, with other genera known from fossils. The extant species fall into three families—the Callorhinchidae, Rhinochimaeridae and Chimaeridae with the callorhinchids being the most basal clade.
Suborder Chimaeroidei Patterson 1965
Family Callorhinchidae Garman, 1901
Genus Callorhinchus Lacépède, 1798 (3 extant species) Mid-Cretaceous–recent
Family Chimaeridae Bonaparte, 1831
Genus Chimaera Linnaeus, 1758 (16 species) Eocene–recent
Genus Hydrolagus Gill, 1863 (26 species) Miocene–recent
Family Rhinochimaeridae Garman, 1901
Genus Harriotta Goode & Bean, 1895 (2 species)
Genus Neoharriotta Bigelow & Schroeder, 1950 (3 species)
Genus Rhinochimaera Garman, 1901 (3 species)
Evolution
Tracing the evolution of these species has been problematic given the paucity of good fossils. DNA sequencing has become the preferred approach to understanding speciation.
The group containing chimaeras and their close relatives (Holocephali) is thought to have diverged from Elasmobranchii (the group containing modern sharks and rays) during the Devonian, over 380 million years ago. The oldest known chimaeriform is Protochimaera from the Early Carboniferous (338–332 million years ago) of Russia, which is more closely related to modern chimeras (Chimaeroidei) than any other known extinct groups of Chimaeriformes. The earliest known remains attributable to modern chimaeras are known from the Early Jurassic (Pliensbachian) of Europe, but egg cases from the Late Triassic of Yakutia, Russia and New Zealand that resemble those of rhinochimaerids and callorhinchids respectively indicates that they had a global distribution prior to the end of the Triassic. Unlike modern chimaeras, Mesozoic representatives are often found in shallow water settings. Most modern chimaera groups appear to have originated during the Mesozoic Marine Revolution. Modern chimaeras reached their highest ecological diversity during the mid-Cretaceous (Albian to Cenomanian), when they acquired a variety of different dentition types.
It has commonly been assumed that due to being an evolutionarily basal group that is largely found in the deep ocean, modern chimaeras likely colonized the deep ocean during the Mesozoic and used it as a refugium to survive mass extinction events. However, more recent studies indicate that chimaeras were likely a shallow-water group for most of their existence, and only colonized the deep ocean in the aftermath of the Cretaceous-Paleogene extinction event. The plough-nosed chimaeras are the only group to still inhabit shallower waters, in the manner of ancestral chimaera groups.
Taxonomy
Extinct chimaeriforms include:
†Suborder Echinochimaeroidei Lund, 1977
†Family Echinochimaeridae Lund, 1977
†Genus Echinochimaera Lund, 1977 United States, Lower Carboniferous (Serpukhovian)
†Suborder Squalorajoidei Patterson, 1965 (Lower Carboniferous–Early Jurassic)
?†Genus Sulcacanthus Itano & Duffin, 2023 United States, Lower Carboniferous (Viséan)
Family †Squalorajidae Woodward, 1886
†Genus Squaloraja Riley, 1833 Europe, Early Jurassic (Hettangian–Sinemurian)
†Suborder Myriacanthoidei Patterson 1965 (Late Triassic–Late Jurassic; possible Carboniferous records)
†Family Chimaeropsidae
†Chimaeropsis Zittel 1887 Belgium, Early Jurassic (Sinemurian)
†Family Myriacanthidae Woodward 1889
†Acanthorhina Fraas 1910 Posidonia Shale Formation, Germany, Early Jurassic (Toarcian)
†Agkistracanthus Duffin and Furrer 1981 Austria, England and Switzerland, Late Triassic–Early Jurassic (Rhaetian–Sinemurian)
†Alethodontus Duffin 1983 Germany, Early Jurassic (Sinemurian)
†Halonodon Duffin 1984 Belgium and Luxembourg, Early Jurassic (Sinemurian)
†Metopacanthus Zittel 1887 Posidonia Shale Formation, Germany, Early Jurassic (Toarcian)
†Oblidens Duffin and Milàn 2017 Hasle Formation, Denmark, Early Jurassic (Pliensbachian)
†Myriacanthus Agassiz 1837 United Kingdom, Late Triassic-Early Jurassic (Rhaetian–Sinemurian)
†Recurvacanthus Duffin 1981 United Kingdom, Early Jurassic (Sinemurian)
†Suborder Protochimaeroidei Lebedev & Popov in Lebedev et al., 2021
†Family Protochimaeridae Lebedev & Popov in Lebedev et al., 2021
†Genus Protochimaera Lebedev & Popov in Lebedev et al., 2021 Moscow Region, Russia, Lower Carboniferous (Viséan–Serpukhovian)
Suborder Chimaeroidei Patterson 1965
†Eomanodon Ward and Duffin 1989 United Kingdom, Early Jurassic (Pliensbachian)
Family Callorhinchidae Garman, 1901
†Brachymylus A. S. Woodward 1894 Germany, Early Jurassic (Pliensbachian)
†Bathytheristes Duffin 1995 Posidonia Shale Formation, Germany, Early Jurassic (Toarcian)
†Ottangodus Popov, Delsate & Felten, 2019 France, Middle Jurassic (Bajocian)
†Moskovirhynchus Russia, Upper Jurassic
†Pachymylus United Kingdom, France, Middle Jurassic
Family †"Edaphodontidae"
†Ischyodus (40 species) Worldwide, Middle Jurassic–Miocene (also placed in Callorhinchidae)
†Elasmodectes Europe, Jurassic–Cretaceous
†Elasmodus Worldwide, Cretaceous–Paleogene
†Edaphodon Worldwide, Cretaceous–Neogene
†Ptyktoptychion Australia, Early Cretaceous
†Lebediodon Europe, Cretaceous
Family Chimaeridae Bonaparte, 1831
†Canadodus Popov, Johns & Suntok, 2020 Sooke Formation, Canada, Oligocene
Family Rhinochimaeridae Garman, 1901
†Amylodon Europe, Late Cretaceous–Oligocene
| Biology and health sciences | Fishes | null |
490442 | https://en.wikipedia.org/wiki/Geometer%20moth | Geometer moth | The geometer moths are moths belonging to the family Geometridae of the insect order Lepidoptera, the moths and butterflies. Their scientific name derives from the Ancient Greek geo γεω (derivative form of or "the earth"), and metron "measure" in reference to the way their larvae, or inchworms, appear to measure the earth as they move along in a looping fashion. Geometridae is a very large family, containing around 23,000 described species; over 1400 species from six subfamilies are indigenous to North America alone. A well-known member is the peppered moth, Biston betularia, which has been the subject of numerous studies in population genetics. Several other geometer moths are notorious pests.
Caterpillars
The name "Geometridae" ultimately derives from Latin from Greek ("geometer", "earth-measurer"). This refers to the means of locomotion of the larvae or caterpillars, which lack the full complement of prolegs seen in other caterpillars, with only two or three pairs at the posterior end instead of the usual five pairs. Equipped with appendages at both ends of the body, a caterpillar clasps with its front legs and draws up the hind end, then clasps with the hind end (prolegs) and reaches out for a new front attachment, creating the impression that it measures its journey. The caterpillars are accordingly called "loopers", "spanworms", or "inchworms" after their characteristic looping gait. The cabbage looper and soybean looper are not inchworms but caterpillars of a different family. In many species of geometer moths, the inchworms are about long. They tend to be green, grey, or brownish and hide from predators by fading into the background or resembling twigs. When disturbed, many inchworms stand erect and motionless on their prolegs, further increasing this resemblance. Some have humps or filaments, or cover themselves in plant material. They are gregarious and are generally smooth. Some eat lichen, flowers, or pollen, while some, such as the Hawaiian species of the genus Eupithecia, are carnivorous. Certain destructive inchworm species are referred to as "cankerworms".
In 2019, the first geometrid caterpillar in Baltic amber was discovered by German scientists. Described under Eogeometer vadens, it measured about and was estimated to be 44 million years old, dating back to the Eocene epoch. It was described as the earliest evidence for the subfamily of Ennominae, particularly the tribe Boarmiini.
Adults
Many geometrids have slender abdomens and broad wings which are usually held flat with the hindwings visible. As such, they appear rather butterfly-like, but in most respects they are typical moths. The majority fly at night. They possess a frenulum to link the wings, and the antennae of the males are often feathered. They tend to blend into the background, often with intricate, wavy patterns on their wings. In some species, females have reduced wings (e.g. winter moth and fall cankerworm). Most are of moderate size, about in wingspan, but a range of sizes occur, from , and a few (e.g., Dysphania species) reach an even larger size. They have distinctive paired tympanal organs at the base of the abdomen (these are absent in flightless females).
Systematics
The placement of the example species follows a 1990 systematic treatment; it may be outdated. Subfamilies are tentatively sorted in a phylogenetic sequence, from the most basal to the most advanced. Traditionally, the Archiearinae were held to be the most ancient of the geometer moth lineages, as their caterpillars have well-developed prolegs. However, it now seems that the Larentiinae are actually older, as indicated by their numerous plesiomorphies and DNA sequence data. They are either an extremely basal lineage of the Geometridae – together with the Sterrhinae – or might even be considered a separate family of Geometroidea. As regards the Archiearinae, some species that were traditionally placed therein actually seem to belong to other subfamilies; altogether it seems that in a few cases, the prolegs which were originally lost in the ancestral geometer moths re-evolved as an atavism.
Larentiinae – about 5,800 species, includes the pug moths, mostly temperate, might be a distinct family.
Sterrhinae – about 2,800 species, mostly tropical, might belong to same family as the Larentiinae.
Birch mocha, Cyclophora albipunctata
False mocha, Cyclophora porata
Maiden's blush, Cyclophora punctaria
Riband wave, Idaea aversata
Small fan-footed wave, Idaea biselata
Single-dotted wave, Idaea dimidiata
Small scallop, Idaea emarginata
Idaea filicata
Dwarf cream wave, Idaea fuscovenosa
Rusty wave, Idaea inquinata
Purple-bordered gold, Idaea muricata
Bright wave, Idaea ochrata
Least carpet, Idaea rusticata
Small dusty wave, Idaea seriata
Purple-barred yellow, Lythria cruentaria (formerly in Larentiinae)
Vestal, Rhodometra sacraria
Common pink-barred, Rhodostrophia vibicaria
Middle lace border, Scopula decorata
Cream wave, Scopula floslactata
Small blood-vein, Scopula imitaria
Lewes wave, Scopula immorata
Lesser cream wave, Scopula immutata
Mullein wave, Scopula marginepunctata
Zachera moth, Chiasmia defixaria
Blood-vein, Timandra comae
Eastern blood-vein, Timandra griseata
Desmobathrinae – pantropical
Geometrinae – emerald moths, about 2,300 named species, most tropical
Archiearinae – twelve species; holarctic, southern Andes and Tasmania, though the latter some seem to belong to the Ennominae, larvae have all the prolegs but most are reduced.
Infant, Archiearis infans (Möschler, 1862)
Scarce infant, Leucobrephos brephoides (Walker, 1857)
Oenochrominae – in some treatments used as a "wastebin taxon" for genera that are difficult to place in other groups
Alsophilinae – a few genera, defoliators of trees, might belong in the Ennominae, tribe Boarmiini
March moth, Alsophila aescularia
Fall cankerworm, Alsophila pometaria
Ennominae – about 9,700 species, including some defoliating pests, global distribution
†Eogeometer vadens
Geometridae genera incertae sedis include:
Dichromodes
Homoeoctenia
Nearcha
Fossil Geometridae taxa include:
†Eogeometer Fischer, Michalski & Hausmann, 2019
†Hydriomena? protrita Cockerell, 1922 (Priabonian, Florissant Formation, Colorado)
†Geometridites Clark et al., 1971
| Biology and health sciences | Lepidoptera | null |
490528 | https://en.wikipedia.org/wiki/VLC%20media%20player | VLC media player | VLC media player (previously the VideoLAN Client and commonly known as simply VLC) is a free and open-source, portable, cross-platform media player software and streaming media server developed by the VideoLAN project. VLC is available for desktop operating systems and mobile platforms, such as Android, iOS and iPadOS. VLC is also available on digital distribution platforms such as Apple's App Store, Google Play, and Microsoft Store.
VLC supports many audio- and video-compression-methods and file-formats, including DVD-Video, Video CD, and streaming-protocols. It is able to stream media over computer networks and can transcode multimedia files.
The default distribution of VLC includes many free decoding and encoding libraries, avoiding the need for finding/calibrating proprietary plugins. The libavcodec library from the FFmpeg project provides many of VLC's codecs, but the player mainly uses its own muxers and demuxers. It also has its own protocol implementations. It also gained distinction as the first player to support playback of encrypted DVDs on Linux and macOS by using the libdvdcss DVD decryption library; however, this library is legally controversial and is not included in many software repositories of Linux distributions as a result. It is available on iOS under the MPLv2.
History
The VideoLAN software originated as a French academic project in 1996. VLC used to stand for "VideoLAN Client" when VLC was a client of the VideoLAN project. Since VLC is no longer merely a client, that initialism no longer applies. It was intended to consist of a client and server to stream videos from satellite dishes across a campus network. Originally developed by students at the École Centrale Paris, it is now developed by contributors worldwide and is coordinated by VideoLAN, a non-profit organization. Rewritten from scratch in 1998, it was released under GNU General Public License on February 1, 2001, with authorization from the headmaster of the École Centrale Paris. The functionality of the server-program, VideoLan Server (VLS), has mostly been subsumed into VLC and has been deprecated. The project name has been changed to VLC media player because there is no longer a client/server infrastructure.
The cone icon used in VLC is a reference to the traffic cones collected by École Centrale's Networking Students' Association. The cone icon design was changed from a hand drawn low resolution icon to a higher resolution CGI-rendered version in 2005, illustrated by Richard Øiestad.
In 2007 the VLC project decided, for license compatibility reasons, not to upgrade to the just-released GPLv3. After 13 years of development, version 1.0.0 of VLC media player was released on July 7, 2009. Work began on VLC for Android in 2010 and it has been available for Android devices on the Google Play store since 2011. In September 2010, a company named "Applidium" developed a VLC port for iOS under GPLv2 with the endorsement of the VLC project, which was accepted by Apple for their App Store. In January 2011, after VLC developer Rémi Denis-Courmont's complaint to Apple about the licensing conflict between the VLC's GPLv2 and the App store's policies, the VLC had been withdrawn from the Apple App Store by Apple. Subsequently, in October 2011 the VLC authors began to relicense the engine parts of VLC from the GPL-2.0-or-later to the LGPL-2.1-or-later to achieve better license compatibility, for instance with the Apple App Store. In July 2013 the VLC application could be resubmitted to the iOS App Store under the MPL-2.0. Version 2.0.0 of VLC media player was released on February 18, 2012. The version for the Windows Store was released on March 13, 2014. Support for Windows RT, Windows Phone and Xbox One were added later. VLC is the third in the sourceforge.net overall download count, and there have been more than 6 billion downloads.
Version 3.0 was in development for Windows, Linux and macOS since June 2016 and released in February 2018. It contains many new features including Chromecast output support (except subtitles), hardware-accelerated decoding enabled by default, 4K and 8K playback, 10-bit and HDR playback, 360° video and 3D audio, audio passthrough for HD audio codecs, BD-J menu support, and local network drive browsing.
In December 2017 the European Parliament approved a budget that funds a bug bounty program for VLC to improve the EU's IT infrastructure.
Release history
Starting with version 1.1.0, VLC release codenames refer to characters from Terry Pratchett's Discworld novels; an exception is release 2.2.1, which came out shortly after Pratchett's death on March 12, 2015, and which was codenamed Terry Pratchett in honor of the author himself.
Design principles
Modular design
VLC, like most multimedia frameworks, has a very modular design which makes it easier to include modules/plugins for new file formats, codecs, interfaces, or streaming methods. VLC 1.0.0 has more than 380 modules. The VLC core creates its own graph of modules dynamically, depending on the situation: input protocol, input file format, input codec, video card capabilities and other parameters. In VLC, almost everything is a module, like interfaces, video and audio outputs, controls, scalers, codecs, and audio/video filters.
Interfaces
The default GUI is based on Be API on BeOS, Cocoa for macOS, and Qt 5 for Linux and Windows, but all give a similar standard interface. The old default GUI was based on wxWidgets on Linux and Windows. VLC supports highly customizable skins through the skins2 interface, and also supports Winamp 2 and XMMS skins. Skins are not supported in the macOS version. VLC has ncurses, remote control, and telnet console interfaces. There is also an HTTP interface, as well as interfaces for mouse gestures and keyboard hotkeys.
Features
Effects (desktop version)
The desktop version of VLC media player has some filters that can distort, rotate, split, deinterlace, and mirror videos as well as create display walls or add a logo overlay during playback. It can also output video as ASCII art.
An interactive zoom feature allows magnifying into video during playback. Still images can be extracted from video at original resolution, and individual frames can be stepped through, although only in forward direction.
Playback can be gamified by splitting the picture inside the viewport into draggable puzzle pieces, where the row and column count can be set as desired.
For audio playback, this feature includes an equalizer and other filters that help customize sound quality.
Formats
Because VLC is a packet-based media player it plays almost all video content. Even some damaged, incomplete, or unfinished files can be played, such as those still downloading via a peer-to-peer (P2P) network. It also plays m2t MPEG transport streams (.TS) files while they are still being digitized from an HDV camera via a FireWire cable, making it possible to monitor the video as it is being recorded. The player can also use libcdio to access .iso files so that users can play files on a disk image, even if the user's operating system cannot work directly with .iso images.
VLC supports all audio and video formats supported by libavcodec and libavformat. This means that VLC can play back H.264 or MPEG-4 Part 2 video as well as support FLV or MXF file formats "out of the box" using FFmpeg's libraries. Alternatively, VLC has modules for codecs that are not based on FFmpeg's libraries. VLC is one of the free software DVD players that ignore DVD region coding on RPC-1 firmware drives, making it a region-free player. However, it does not do the same on RPC-2 firmware drives, as in these cases the region coding is enforced by the drive itself, however, it can still brute-force the CSS encryption to play a foreign-region DVD on an RPC-2 drive.
VLC media player can play high-definition recordings of D-VHS tapes duplicated to a computer using . This offers another way to archive all D-VHS tapes with the DRM copy freely tag. Using a FireWire connection from cable boxes to computers, VLC can stream live, unencrypted content to a monitor or HDTV. VLC media player can display the playing video as the desktop wallpaper, like Windows DreamScene, by using DirectX, only available on Windows operating systems. VLC media player can record the desktop and save the stream as a file, allowing the user to create screencasts. On Microsoft Windows, VLC also supports the Direct Media Object (DMO) framework and can thus make use of some third-party DLLs (Dynamic-link library). On most platforms, VLC can tune into and view DVB-C, DVB-T, and DVB-S channels. On macOS the separate EyeTV plugin is required, on Windows it requires the card's BDA Drivers.
VLC can be installed or run directly from a USB flash drive or other external drive. VLC can be extended through scripting; it uses the Lua scripting language. VLC can play videos in the AVCHD format, a highly compressed format used in recent HD camcorders. VLC can generate a number of music visualization displays. The program is able to convert media files into various supported formats.
Both desktop and mobile releases are equipped with an audio equalizer.
Christmas logo
A red Santa hat appears on top of VLC's traffic-cone logo during Christmas seasons.
Keyboard shortcuts
There are single-button shortcuts in VLC that don't require Ctrl or Alt button.
For example, pressing keys F and G while a video file is running in VLC shifts the file's audio/video sync for 50 millisecond per adjustment. This is useful to fix an issue with the sound being ahead or lagging behind the video.
Operating system compatibility
VLC media player is cross-platform, with versions for Windows, macOS, Linux, iOS, Android, tvOS, ChromeOS, Windows Phone, various BSD-based systems, Solaris, BeOS, OS/2, and Syllable. However, forward and backward compatibility between versions of VLC media player and different versions of OSes are not maintained over more than a few generations. 64-bit builds are available for 64-bit Windows, starting with version 2.0.1.
Windows 8 and 10 support
The VLC port for Windows 8 and Windows 10 is backed by a crowdfunding campaign on Kickstarter to add support for a new GUI based on Microsoft's Metro design language, that will run on the Windows Runtime. All the existing features including video filters, subtitle support, and an equalizer are present in Windows 8. A beta version of VLC for Windows 8 was released to the Microsoft Store on March 13, 2014. A universal app was created for Windows 8, 8.1, 10, Windows Phone 8, 8.1 and Windows 10 Mobile.
Android support
In May 2012, the VLC team stated that a version of VLC for Android was being developed. The stable release version 1.0 was made available on Google Play on December 8, 2014.
Use of VLC with other programs
Bindings
Several APIs can connect to VLC and use its functionality:
libVLC API – the VLC Core, for C and C++
VLCKit – an Objective-C framework for macOS
LibVLCSharp – Crossplatform .NET bindings to libVLC (C#/F#/VB)
JavaScript API – the evolution of ActiveX API and Firefox integration
D-Bus controls
Go bindings
Python controls
Java API
DirectShow filters
Delphi/Pascal API: PasLibVlc by Robert Jędrzejczyk
Free Pascal bindings and an OOP wrapper component, via the libvlc.pp and vlc.pp units. This comes standard with the Free Pascal Compiler as of November 6, 2012.
The Phonon multimedia API for Qt and KDE applications can optionally use VLC as a backend.
Applications that use libVLC
VLC can handle some incomplete files and in some cases can be used to preview files being downloaded. Several programs make use of this, including eMule and KCeasy. The free/open-source Internet television application Miro also uses VLC code. HandBrake, an open-source video encoder, used to load libdvdcss from VLC Media Player. Easy Subtitles Synchronizer, a freeware subtitle editing program for Windows, uses VLC to preview the video with the edited subtitles.
Format support
Input formats
VLC can read many formats, depending on the operating system it is running on, including:
Container formats: 3GP, ASF, AVI, DVR-MS, FLV, Matroska (MKV), MIDI, QuickTime File Format, MP4, Ogg, OGM, WAV, MPEG-2 (ES, PS, TS, PVA, MP3), AIFF, Raw audio, Raw DV, MXF, VOB, RM, Blu-ray, DVD-Video, VCD, SVCD, CD-DA, DVB, HEIF, AVIF
Audio coding formats: AAC, AC3, ALAC, AMR, DTS, DV Audio, XM, FLAC, It, MACE, MOD, Monkey's Audio, MP3, Opus, PLS, QCP, QDM2/QDMC, RealAudio, Speex, Screamtracker 3/S3M, TTA, Vorbis, WavPack, WMA (WMA 1/2, WMA 3 partially).
Capture devices: Video4Linux (on Linux), DirectShow (on Windows), Desktop (screencast), Digital TV (DVB-C, DVB-S, DVB-T, DVB-S2, DVB-T2, ATSC, Clear QAM)
Network protocols: FTP, HTTP, MMS, RSS/Atom, RTMP, RTP (unicast or multicast), RTSP, UDP, Sat-IP, Smooth Streaming
Network streaming formats: Apple HLS, Flash RTMP, MPEG-DASH, MPEG Transport Stream, RTP/RTSP ISMA/3GPP PSS, Windows Media MMS
Subtitles: Advanced SubStation Alpha, Closed Captions, DVB, DVD-Video, MPEG-4 Timed Text, MPL2, OGM, SubStation Alpha, SubRip, SVCD, Teletext, Text file, VobSub, WebVTT, TTML
Video coding formats: Cinepak, Dirac, DV, H.263, H.264/MPEG-4 AVC, H.265/MPEG HEVC, AV1, HuffYUV, Indeo 3, MJPEG, MPEG-1, MPEG-2, MPEG-4 Part 2, RealVideo 3&4, Sorenson, Theora, VC-1, VP5, VP6, VP8, VP9, DNxHD, ProRes and some WMV.
Digital Camcorder formats: MOD and TOD via USB.
Output formats
VLC can transcode or stream audio and video into several formats depending on the operating system, including:
Container formats: ASF, AVI, FLAC, FLV, Fraps, Matroska, MP4, MPJPEG, MPEG-2 (ES, MP3), Ogg, PS, PVA, QuickTime File Format, TS, WAV, WebM
Audio coding formats: AAC, AC-3, DV Audio, FLAC, MP3, Speex, Vorbis
Streaming protocols: HTTP, MMS, RTSP, RTP, UDP
Video coding formats: Dirac, DV, H.263, H.264/MPEG-4 AVC, H.265/MPEG-H HEVC, MJPEG, MPEG-1, MPEG-2, MPEG-4 Part 2, Theora, VP5, VP6, VP8, VP9
Legality
The VLC media player software installers for the macOS platform and the Windows platform include the libdvdcss DVD decryption library, even though this library may be legally restricted in certain jurisdictions.
India
In May 2022, it was reported by MediaNama that VLC was banned in India and its website was inaccessible from India under the provisions of the Information Technology Act, 2000. Neither the developers nor the Indian government offered any explanation to the ban, according to India Today. The official VideoLAN Twitter account stated in August that the website was blocked in India from 13 February 2022. A report by Hindustan Times indicated that the ban could be due to links with China. India had in 2020 banned over 200 Chinese apps following the 2020–2022 China–India skirmishes.
Another Hindustan Times report from April quoting Symantec said that Chinese hackers were depending on VLC to launch malware they had previously installed on Windows machines. The technique they used is called DLL side-loading, in which an external library that a legitimate program loads at runtime is substituted with a modified version containing the malware.
VideoLan president and lead developer Jean-Baptiste Kempf said that the block was most likely a result of a misunderstanding of the Chinese security issue, although the Indian Government did not provide for a reason as to why it was blocked. In October 2022, VideoLan, with assistance from the Indian digital rights organization Internet Freedom Foundation sent a legal notice to the Indian government asking for an explanation for the block order, following which the Ministry of Electronics and Information Technology removed the ban in November 2022.
United States
The VLC media player software is able to read audio and video data from DVDs that incorporate Content Scramble System (CSS) encryption, even though the VLC media player software lacks a CSS decryption license. The unauthorized decryption of CSS-encrypted DVD content or unauthorized distribution of CSS decryption tools may violate the US Digital Millennium Copyright Act.
Decryption of CSS-encrypted DVD content has been temporarily authorized for certain purposes (such as documentary filmmaking that uses short portions of DVD content for criticism or commentary) under the Digital Millennium Copyright Act anticircumvention exemptions that were issued by the US Copyright Office in 2010. However, these exemptions do not change the DMCA's ban on the distribution of CSS decryption tools; including those distributed with VLC.
| Technology | Multimedia_2 | null |
490620 | https://en.wikipedia.org/wiki/Human%20brain | Human brain | The human brain is the central organ of the human nervous system, and with the spinal cord, comprises the central nervous system. It consists of the cerebrum, the brainstem and the cerebellum. The brain controls most of the activities of the body, processing, integrating, and coordinating the information it receives from the sensory nervous system. The brain integrates the instructions sent to the rest of the body. The brain is contained in, and protected by, the skull of the head.
The cerebrum, the largest part of the human brain, consists of two cerebral hemispheres. Each hemisphere has an inner core composed of white matter, and an outer surface – the cerebral cortex – composed of grey matter. The cortex has an outer layer, the neocortex, and an inner allocortex. The neocortex is made up of six neuronal layers, while the allocortex has three or four. Each hemisphere is divided into four lobes – the frontal, parietal, temporal, and occipital lobes. The frontal lobe is associated with executive functions including self-control, planning, reasoning, and abstract thought, while the occipital lobe is dedicated to vision. Within each lobe, cortical areas are associated with specific functions, such as the sensory, motor, and association regions. Although the left and right hemispheres are broadly similar in shape and function, some functions are associated with one side, such as language in the left and visual-spatial ability in the right. The hemispheres are connected by commissural nerve tracts, the largest being the corpus callosum.
The cerebrum is connected by the brainstem to the spinal cord. The brainstem consists of the midbrain, the pons, and the medulla oblongata. The cerebellum is connected to the brainstem by three pairs of nerve tracts called cerebellar peduncles. Within the cerebrum is the ventricular system, consisting of four interconnected ventricles in which cerebrospinal fluid is produced and circulated. Underneath the cerebral cortex are several structures, including the thalamus, the epithalamus, the pineal gland, the hypothalamus, the pituitary gland, and the subthalamus; the limbic structures, including the amygdalae and the hippocampi, the claustrum, the various nuclei of the basal ganglia, the basal forebrain structures, and three circumventricular organs. Brain structures that are not on the midplane exist in pairs; for example, there are two hippocampi and two amygdalae.
The cells of the brain include neurons and supportive glial cells. There are more than 86 billion neurons in the brain, and a more or less equal number of other cells. Brain activity is made possible by the interconnections of neurons and their release of neurotransmitters in response to nerve impulses. Neurons connect to form neural pathways, neural circuits, and elaborate network systems. The whole circuitry is driven by the process of neurotransmission.
The brain is protected by the skull, suspended in cerebrospinal fluid, and isolated from the bloodstream by the blood–brain barrier. However, the brain is still susceptible to damage, disease, and infection. Damage can be caused by trauma, or a loss of blood supply known as a stroke. The brain is susceptible to degenerative disorders, such as Parkinson's disease, dementias including Alzheimer's disease, and multiple sclerosis. Psychiatric conditions, including schizophrenia and clinical depression, are thought to be associated with brain dysfunctions. The brain can also be the site of tumours, both benign and malignant; these mostly originate from other sites in the body.
The study of the anatomy of the brain is neuroanatomy, while the study of its function is neuroscience. Numerous techniques are used to study the brain. Specimens from other animals, which may be examined microscopically, have traditionally provided much information. Medical imaging technologies such as functional neuroimaging, and electroencephalography (EEG) recordings are important in studying the brain. The medical history of people with brain injury has provided insight into the function of each part of the brain. Neuroscience research has expanded considerably, and research is ongoing.
In culture, the philosophy of mind has for centuries attempted to address the question of the nature of consciousness and the mind–body problem. The pseudoscience of phrenology attempted to localise personality attributes to regions of the cortex in the 19th century. In science fiction, brain transplants are imagined in tales such as the 1942 Donovan's Brain.
Structure
Gross anatomy
The adult human brain weighs on average about which is about 2% of the total body weight, with a volume of around 1260 cm3 in men and 1130 cm3 in women. There is substantial individual variation, with the standard reference range for men being and for women .
The cerebrum, consisting of the cerebral hemispheres, forms the largest part of the brain and overlies the other brain structures. The outer region of the hemispheres, the cerebral cortex, is grey matter, consisting of cortical layers of neurons. Each hemisphere is divided into four main lobes – the frontal lobe, parietal lobe, temporal lobe, and occipital lobe. Three other lobes are included by some sources which are a central lobe, a limbic lobe, and an insular lobe. The central lobe comprises the precentral gyrus and the postcentral gyrus and is included since it forms a distinct functional role.
The brainstem, resembling a stalk, attaches to and leaves the cerebrum at the start of the midbrain area. The brainstem includes the midbrain, the pons, and the medulla oblongata. Behind the brainstem is the cerebellum ().
The cerebrum, brainstem, cerebellum, and spinal cord are covered by three membranes called meninges. The membranes are the tough dura mater; the middle arachnoid mater and the more delicate inner pia mater. Between the arachnoid mater and the pia mater is the subarachnoid space and subarachnoid cisterns, which contain the cerebrospinal fluid. The outermost membrane of the cerebral cortex is the basement membrane of the pia mater called the glia limitans and is an important part of the blood–brain barrier. In 2023 a fourth meningeal membrane has been proposed known as the subarachnoid lymphatic-like membrane. The living brain is very soft, having a gel-like consistency similar to soft tofu. The cortical layers of neurons constitute much of the cerebral grey matter, while the deeper subcortical regions of myelinated axons, make up the white matter. The white matter of the brain makes up about half of the total brain volume.
Cerebrum
The cerebrum is the largest part of the brain and is divided into nearly symmetrical left and right hemispheres by a deep groove, the longitudinal fissure. Asymmetry between the lobes is noted as a petalia. The hemispheres are connected by five commissures that span the longitudinal fissure, the largest of these is the corpus callosum.
Each hemisphere is conventionally divided into four main lobes; the frontal lobe, parietal lobe, temporal lobe, and occipital lobe, named according to the skull bones that overlie them. Each lobe is associated with one or two specialised functions though there is some functional overlap between them. The surface of the brain is folded into ridges (gyri) and grooves (sulci), many of which are named, usually according to their position, such as the frontal gyrus of the frontal lobe or the central sulcus separating the central regions of the hemispheres. There are many small variations in the secondary and tertiary folds.
The outer part of the cerebrum is the cerebral cortex, made up of grey matter arranged in layers. It is thick, and deeply folded to give a convoluted appearance. Beneath the cortex is the cerebral white matter. The largest part of the cerebral cortex is the neocortex, which has six neuronal layers. The rest of the cortex is of allocortex, which has three or four layers.
The cortex is mapped by divisions into about fifty different functional areas known as Brodmann's areas. These areas are distinctly different when seen under a microscope. The cortex is divided into two main functional areas – a motor cortex and a sensory cortex. The primary motor cortex, which sends axons down to motor neurons in the brainstem and spinal cord, occupies the rear portion of the frontal lobe, directly in front of the somatosensory area. The primary sensory areas receive signals from the sensory nerves and tracts by way of relay nuclei in the thalamus. Primary sensory areas include the visual cortex of the occipital lobe, the auditory cortex in parts of the temporal lobe and insular cortex, and the somatosensory cortex in the parietal lobe. The remaining parts of the cortex are called the association areas. These areas receive input from the sensory areas and lower parts of the brain and are involved in the complex cognitive processes of perception, thought, and decision-making. The main functions of the frontal lobe are to control attention, abstract thinking, behaviour, problem-solving tasks, and physical reactions and personality. The occipital lobe is the smallest lobe; its main functions are visual reception, visual-spatial processing, movement, and colour recognition. There is a smaller occipital lobule in the lobe known as the cuneus. The temporal lobe controls auditory and visual memories, language, and some hearing and speech.
The cerebrum contains the ventricles where the cerebrospinal fluid is produced and circulated. Below the corpus callosum is the septum pellucidum, a membrane that separates the lateral ventricles. Beneath the lateral ventricles is the thalamus and to the front and below is the hypothalamus. The hypothalamus leads on to the pituitary gland. At the back of the thalamus is the brainstem.
The basal ganglia, also called basal nuclei, are a set of structures deep within the hemispheres involved in behaviour and movement regulation. The largest component is the striatum, others are the globus pallidus, the substantia nigra and the subthalamic nucleus. The striatum is divided into a ventral striatum, and dorsal striatum, subdivisions that are based upon function and connections. The ventral striatum consists of the nucleus accumbens and the olfactory tubercle whereas the dorsal striatum consists of the caudate nucleus and the putamen. The putamen and the globus pallidus lie separated from the lateral ventricles and thalamus by the internal capsule, whereas the caudate nucleus stretches around and abuts the lateral ventricles on their outer sides. At the deepest part of the lateral sulcus between the insular cortex and the striatum is a thin neuronal sheet called the claustrum.
Below and in front of the striatum are a number of basal forebrain structures. These include the nucleus basalis, diagonal band of Broca, substantia innominata, and the medial septal nucleus. These structures are important in producing the neurotransmitter, acetylcholine, which is then distributed widely throughout the brain. The basal forebrain, in particular the nucleus basalis, is considered to be the major cholinergic output of the central nervous system to the striatum and neocortex.
Cerebellum
The cerebellum is divided into an anterior lobe, a posterior lobe, and the flocculonodular lobe. The anterior and posterior lobes are connected in the middle by the vermis. Compared to the cerebral cortex, the cerebellum has a much thinner outer cortex that is narrowly furrowed into numerous curved transverse fissures.
Viewed from underneath between the two lobes is the third lobe the flocculonodular lobe. The cerebellum rests at the back of the cranial cavity, lying beneath the occipital lobes, and is separated from these by the cerebellar tentorium, a sheet of fibre.
It is connected to the brainstem by three pairs of nerve tracts called cerebellar peduncles. The superior pair connects to the midbrain; the middle pair connects to the medulla, and the inferior pair connects to the pons. The cerebellum consists of an inner medulla of white matter and an outer cortex of richly folded grey matter. The cerebellum's anterior and posterior lobes appear to play a role in the coordination and smoothing of complex motor movements, and the flocculonodular lobe in the maintenance of balance although debate exists as to its cognitive, behavioural and motor functions.
Brainstem
The brainstem lies beneath the cerebrum and consists of the midbrain, pons and medulla. It lies in the back part of the skull, resting on the part of the base known as the clivus, and ends at the foramen magnum, a large opening in the occipital bone. The brainstem continues below this as the spinal cord, protected by the vertebral column.
Ten of the twelve pairs of cranial nerves emerge directly from the brainstem. The brainstem also contains many cranial nerve nuclei and nuclei of peripheral nerves, as well as nuclei involved in the regulation of many essential processes including breathing, control of eye movements and balance. The reticular formation, a network of nuclei of ill-defined formation, is present within and along the length of the brainstem. Many nerve tracts, which transmit information to and from the cerebral cortex to the rest of the body, pass through the brainstem.
Microanatomy
The human brain is primarily composed of neurons, glial cells, neural stem cells, and blood vessels. Types of neuron include interneurons, pyramidal cells including Betz cells, motor neurons (upper and lower motor neurons), and cerebellar Purkinje cells. Betz cells are the largest cells (by size of cell body) in the nervous system. The adult human brain is estimated to contain 86±8 billion neurons, with a roughly equal number (85±10 billion) of non-neuronal cells. Out of these neurons, 16 billion (19%) are located in the cerebral cortex, and 69 billion (80%) are in the cerebellum.
Types of glial cell are astrocytes (including Bergmann glia), oligodendrocytes, ependymal cells (including tanycytes), radial glial cells, microglia, and a subtype of oligodendrocyte progenitor cells. Astrocytes are the largest of the glial cells. They are stellate cells with many processes radiating from their cell bodies. Some of these processes end as perivascular endfeet on capillary walls. The glia limitans of the cortex is made up of astrocyte endfeet processes that serve in part to contain the cells of the brain.
Mast cells are white blood cells that interact in the neuroimmune system in the brain. Mast cells in the central nervous system are present in a number of structures including the meninges; they mediate neuroimmune responses in inflammatory conditions and help to maintain the blood–brain barrier, particularly in brain regions where the barrier is absent. Mast cells serve the same general functions in the body and central nervous system, such as effecting or regulating allergic responses, innate and adaptive immunity, autoimmunity, and inflammation. Mast cells serve as the main effector cell through which pathogens can affect the biochemical signaling that takes place between the gastrointestinal tract and the central nervous system.
Some 400 genes are shown to be brain-specific. In all neurons, ELAVL3 is expressed, and in pyramidal cells, NRGN and REEP2 are also expressed. GAD1 – essential for the biosynthesis of the neurotransmitter GABA – is expressed in interneurons. Proteins expressed in glial cells include astrocyte markers GFAP and S100B whereas myelin basic protein and the transcription factor OLIG2 are expressed in oligodendrocytes.
Cerebrospinal fluid
Cerebrospinal fluid is a clear, colourless transcellular fluid that circulates around the brain in the subarachnoid space, in the ventricular system, and in the central canal of the spinal cord. It also fills some gaps in the subarachnoid space, known as subarachnoid cisterns. The four ventricles, two lateral, a third, and a fourth ventricle, all contain a choroid plexus that produces cerebrospinal fluid. The third ventricle lies in the midline and is connected to the lateral ventricles. A single duct, the cerebral aqueduct between the pons and the cerebellum, connects the third ventricle to the fourth ventricle. Three separate openings, the middle and two lateral apertures, drain the cerebrospinal fluid from the fourth ventricle to the cisterna magna, one of the major cisterns. From here, cerebrospinal fluid circulates around the brain and spinal cord in the subarachnoid space, between the arachnoid mater and pia mater.
At any one time, there is about 150mL of cerebrospinal fluid – most within the subarachnoid space. It is constantly being regenerated and absorbed, and is replaced about once every 5–6 hours.
A glymphatic system has been described as the lymphatic drainage system of the brain. The brain-wide glymphatic pathway includes drainage routes from the cerebrospinal fluid, and from the meningeal lymphatic vessels that are associated with the dural sinuses, and run alongside the cerebral blood vessels. The pathway drains interstitial fluid from the tissue of the brain.
Blood supply
The internal carotid arteries supply oxygenated blood to the front of the brain and the vertebral arteries supply blood to the back of the brain. These two circulations join in the circle of Willis, a ring of connected arteries that lies in the interpeduncular cistern between the midbrain and pons.
The internal carotid arteries are branches of the common carotid arteries. They enter the cranium through the carotid canal, travel through the cavernous sinus and enter the subarachnoid space. They then enter the circle of Willis, with two branches, the anterior cerebral arteries emerging. These branches travel forward and then upward along the longitudinal fissure, and supply the front and midline parts of the brain. One or more small anterior communicating arteries join the two anterior cerebral arteries shortly after they emerge as branches. The internal carotid arteries continue forward as the middle cerebral arteries. They travel sideways along the sphenoid bone of the eye socket, then upwards through the insula cortex, where final branches arise. The middle cerebral arteries send branches along their length.
The vertebral arteries emerge as branches of the left and right subclavian arteries. They travel upward through transverse foramina which are spaces in the cervical vertebrae. Each side enters the cranial cavity through the foramen magnum along the corresponding side of the medulla. They give off one of the three cerebellar branches. The vertebral arteries join in front of the middle part of the medulla to form the larger basilar artery, which sends multiple branches to supply the medulla and pons, and the two other anterior and superior cerebellar branches. Finally, the basilar artery divides into two posterior cerebral arteries. These travel outwards, around the superior cerebellar peduncles, and along the top of the cerebellar tentorium, where it sends branches to supply the temporal and occipital lobes. Each posterior cerebral artery sends a small posterior communicating artery to join with the internal carotid arteries.
Blood drainage
Cerebral veins drain deoxygenated blood from the brain. The brain has two main networks of veins: an exterior or superficial network, on the surface of the cerebrum that has three branches, and an interior network. These two networks communicate via anastomosing (joining) veins. The veins of the brain drain into larger cavities of the dural venous sinuses usually situated between the dura mater and the covering of the skull. Blood from the cerebellum and midbrain drains into the great cerebral vein. Blood from the medulla and pons of the brainstem have a variable pattern of drainage, either into the spinal veins or into adjacent cerebral veins.
The blood in the deep part of the brain drains, through a venous plexus into the cavernous sinus at the front, and the superior and inferior petrosal sinuses at the sides, and the inferior sagittal sinus at the back. Blood drains from the outer brain into the large superior sagittal sinus, which rests in the midline on top of the brain. Blood from here joins with blood from the straight sinus at the confluence of sinuses.
Blood from here drains into the left and right transverse sinuses. These then drain into the sigmoid sinuses, which receive blood from the cavernous sinus and superior and inferior petrosal sinuses. The sigmoid drains into the large internal jugular veins.
The blood–brain barrier
The larger arteries throughout the brain supply blood to smaller capillaries. These smallest of blood vessels in the brain, are lined with cells joined by tight junctions and so fluids do not seep in or leak out to the same degree as they do in other capillaries; this creates the blood–brain barrier. Pericytes play a major role in the formation of the tight junctions. The barrier is less permeable to larger molecules, but is still permeable to water, carbon dioxide, oxygen, and most fat-soluble substances (including anaesthetics and alcohol). The blood-brain barrier is not present in the circumventricular organswhich are structures in the brain that may need to respond to changes in body fluidssuch as the pineal gland, area postrema, and some areas of the hypothalamus. There is a similar blood–cerebrospinal fluid barrier, which serves the same purpose as the blood–brain barrier, but facilitates the transport of different substances into the brain due to the distinct structural characteristics between the two barrier systems.
Development
At the beginning of the third week of development, the embryonic ectoderm forms a thickened strip called the neural plate. By the fourth week of development the neural plate has widened to give a broad cephalic end, a less broad middle part and a narrow caudal end. These swellings are known as the primary brain vesicles and represent the beginnings of the forebrain (prosencephalon), midbrain (mesencephalon), and hindbrain (rhombencephalon).
Neural crest cells (derived from the ectoderm) populate the lateral edges of the plate at the neural folds. In the fourth weekduring the neurulation stagethe neural folds close to form the neural tube, bringing together the neural crest cells at the neural crest. The neural crest runs the length of the tube with cranial neural crest cells at the cephalic end and caudal neural crest cells at the tail. Cells detach from the crest and migrate in a craniocaudal (head to tail) wave inside the tube. Cells at the cephalic end give rise to the brain, and cells at the caudal end give rise to the spinal cord.
The tube flexes as it grows, forming the crescent-shaped cerebral hemispheres at the head. The cerebral hemispheres first appear on day 32.
Early in the fourth week, the cephalic part bends sharply forward in a cephalic flexure. This flexed part becomes the forebrain (prosencephalon); the adjoining curving part becomes the midbrain (mesencephalon) and the part caudal to the flexure becomes the hindbrain (rhombencephalon). These areas are formed as swellings known as the three primary brain vesicles. In the fifth week of development five secondary brain vesicles have formed. The forebrain separates into two vesicles – an anterior telencephalon and a posterior diencephalon. The telencephalon gives rise to the cerebral cortex, basal ganglia, and related structures. The diencephalon gives rise to the thalamus and hypothalamus. The hindbrain also splits into two areas – the metencephalon and the myelencephalon. The metencephalon gives rise to the cerebellum and pons. The myelencephalon gives rise to the medulla oblongata. Also during the fifth week, the brain divides into repeating segments called neuromeres. In the hindbrain these are known as rhombomeres.
A characteristic of the brain is the cortical folding known as gyrification. For just over five months of prenatal development the cortex is smooth. By the gestational age of 24 weeks, the wrinkled morphology showing the fissures that begin to mark out the lobes of the brain is evident. Why the cortex wrinkles and folds is not well-understood, but gyrification has been linked to intelligence and neurological disorders, and a number of gyrification theories have been proposed. These theories include those based on mechanical buckling, axonal tension, and differential tangential expansion. What is clear is that gyrification is not a random process, but rather a complex developmentally predetermined process which generates patterns of folds that are consistent between individuals and most species.
The first groove to appear in the fourth month is the lateral cerebral fossa. The expanding caudal end of the hemisphere has to curve over in a forward direction to fit into the restricted space. This covers the fossa and turns it into a much deeper ridge known as the lateral sulcus and this marks out the temporal lobe. By the sixth month other sulci have formed that demarcate the frontal, parietal, and occipital lobes. A gene present in the human genome (ARHGAP11B) may play a major role in gyrification and encephalisation.
Function
Motor control
The frontal lobe is involved in reasoning, motor control, emotion, and language. It contains the motor cortex, which is involved in planning and coordinating movement; the prefrontal cortex, which is responsible for higher-level cognitive functioning; and Broca's area, which is essential for language production. The motor system of the brain is responsible for the generation and control of movement. Generated movements pass from the brain through nerves to motor neurons in the body, which control the action of muscles. The corticospinal tract carries movements from the brain, through the spinal cord, to the torso and limbs. The cranial nerves carry movements related to the eyes, mouth and face.
Gross movement – such as locomotion and the movement of arms and legs – is generated in the motor cortex, divided into three parts: the primary motor cortex, found in the precentral gyrus and has sections dedicated to the movement of different body parts. These movements are supported and regulated by two other areas, lying anterior to the primary motor cortex: the premotor area and the supplementary motor area. The hands and mouth have a much larger area dedicated to them than other body parts, allowing finer movement; this has been visualised in a motor homunculus. Impulses generated from the motor cortex travel along the corticospinal tract along the front of the medulla and cross over (decussate) at the medullary pyramids. These then travel down the spinal cord, with most connecting to interneurons, in turn connecting to lower motor neurons within the grey matter that then transmit the impulse to move to muscles themselves. The cerebellum and basal ganglia, play a role in fine, complex and coordinated muscle movements. Connections between the cortex and the basal ganglia control muscle tone, posture and movement initiation, and are referred to as the extrapyramidal system.
Sensory
The sensory nervous system is involved with the reception and processing of sensory information. This information is received through the cranial nerves, through tracts in the spinal cord, and directly at centres of the brain exposed to the blood. The brain also receives and interprets information from the special senses of vision, smell, hearing, and taste. Mixed motor and sensory signals are also integrated.
From the skin, the brain receives information about fine touch, pressure, pain, vibration and temperature. From the joints, the brain receives information about joint position. The sensory cortex is found just near the motor cortex, and, like the motor cortex, has areas related to sensation from different body parts. Sensation collected by a sensory receptor on the skin is changed to a nerve signal, that is passed up a series of neurons through tracts in the spinal cord. The dorsal column–medial lemniscus pathway contains information about fine touch, vibration and position of joints. The pathway fibres travel up the back part of the spinal cord to the back part of the medulla, where they connect with second-order neurons that immediately send fibres across the midline. These fibres then travel upwards into the ventrobasal complex in the thalamus where they connect with third-order neurons which send fibres up to the sensory cortex. The spinothalamic tract carries information about pain, temperature, and gross touch. The pathway fibres travel up the spinal cord and connect with second-order neurons in the reticular formation of the brainstem for pain and temperature, and also terminate at the ventrobasal complex of the thalamus for gross touch.
Vision is generated by light that hits the retina of the eye. Photoreceptors in the retina transduce the sensory stimulus of light into an electrical nerve signal that is sent to the visual cortex in the occipital lobe. Visual signals leave the retinas through the optic nerves.
Optic nerve fibres from the retinas' nasal halves cross to the opposite sides joining the fibres from the temporal halves of the opposite retinas to form the optic tracts.
The arrangements of the eyes' optics and the visual pathways mean vision from the left visual field is received by the right half of each retina, is processed by the right visual cortex, and vice versa.
The optic tract fibres reach the brain at the lateral geniculate nucleus, and travel through the optic radiation to reach the visual cortex.
Hearing and balance are both generated in the inner ear. Sound results in vibrations of the ossicles which continue finally to the hearing organ, and change in balance results in movement of liquids within the inner ear. This creates a nerve signal that passes through the vestibulocochlear nerve. From here, it passes through to the cochlear nuclei, the superior olivary nucleus, the medial geniculate nucleus, and finally the auditory radiation to the auditory cortex.
The sense of smell is generated by receptor cells in the epithelium of the olfactory mucosa in the nasal cavity. This information passes via the olfactory nerve which goes into the skull through a relatively permeable part. This nerve transmits to the neural circuitry of the olfactory bulb from where information is passed to the olfactory cortex.
Taste is generated from receptors on the tongue and passed along the facial and glossopharyngeal nerves into the solitary nucleus in the brainstem. Some taste information is also passed from the pharynx into this area via the vagus nerve. Information is then passed from here through the thalamus into the gustatory cortex.
Regulation
Autonomic functions of the brain include the regulation, or rhythmic control of the heart rate and rate of breathing, and maintaining homeostasis.
Blood pressure and heart rate are influenced by the vasomotor centre of the medulla, which causes arteries and veins to be somewhat constricted at rest. It does this by influencing the sympathetic and parasympathetic nervous systems via the vagus nerve. Information about blood pressure is generated by baroreceptors in aortic bodies in the aortic arch, and passed to the brain along the afferent fibres of the vagus nerve. Information about the pressure changes in the carotid sinus comes from carotid bodies located near the carotid artery and this is passed via a nerve joining with the glossopharyngeal nerve. This information travels up to the solitary nucleus in the medulla. Signals from here influence the vasomotor centre to adjust vein and artery constriction accordingly.
The brain controls the rate of breathing, mainly by respiratory centres in the medulla and pons. The respiratory centres control respiration, by generating motor signals that are passed down the spinal cord, along the phrenic nerve to the diaphragm and other muscles of respiration. This is a mixed nerve that carries sensory information back to the centres. There are four respiratory centres, three with a more clearly defined function, and an apneustic centre with a less clear function. In the medulla a dorsal respiratory group causes the desire to breathe in and receives sensory information directly from the body. Also in the medulla, the ventral respiratory group influences breathing out during exertion. In the pons the pneumotaxic centre influences the duration of each breath, and the apneustic centre seems to have an influence on inhalation. The respiratory centres directly senses blood carbon dioxide and pH. Information about blood oxygen, carbon dioxide and pH levels are also sensed on the walls of arteries in the peripheral chemoreceptors of the aortic and carotid bodies. This information is passed via the vagus and glossopharyngeal nerves to the respiratory centres. High carbon dioxide, an acidic pH, or low oxygen stimulate the respiratory centres. The desire to breathe in is also affected by pulmonary stretch receptors in the lungs which, when activated, prevent the lungs from overinflating by transmitting information to the respiratory centres via the vagus nerve.
The hypothalamus in the diencephalon, is involved in regulating many functions of the body. Functions include neuroendocrine regulation, regulation of the circadian rhythm, control of the autonomic nervous system, and the regulation of fluid, and food intake. The circadian rhythm is controlled by two main cell groups in the hypothalamus. The anterior hypothalamus includes the suprachiasmatic nucleus and the ventrolateral preoptic nucleus which through gene expression cycles, generates a roughly 24 hour circadian clock. In the circadian day an ultradian rhythm takes control of the sleeping pattern. Sleep is an essential requirement for the body and brain and allows the closing down and resting of the body's systems. There are also findings that suggest that the daily build-up of toxins in the brain are removed during sleep. Whilst awake the brain consumes a fifth of the body's total energy needs. Sleep necessarily reduces this use and gives time for the restoration of energy-giving ATP. The effects of sleep deprivation show the absolute need for sleep.
The lateral hypothalamus contains orexinergic neurons that control appetite and arousal through their projections to the ascending reticular activating system. The hypothalamus controls the pituitary gland through the release of peptides such as oxytocin, and vasopressin, as well as dopamine into the median eminence. Through the autonomic projections, the hypothalamus is involved in regulating functions such as blood pressure, heart rate, breathing, sweating, and other homeostatic mechanisms. The hypothalamus also plays a role in thermal regulation, and when stimulated by the immune system, is capable of generating a fever. The hypothalamus is influenced by the kidneys: when blood pressure falls, the renin released by the kidneys stimulates a need to drink. The hypothalamus also regulates food intake through autonomic signals, and hormone release by the digestive system.
Language
While language functions were traditionally thought to be localised to Wernicke's area and Broca's area, it is now mostly accepted that a wider network of cortical regions contributes to language functions.
The study on how language is represented, processed, and acquired by the brain is called neurolinguistics, which is a large multidisciplinary field drawing from cognitive neuroscience, cognitive linguistics, and psycholinguistics.
Lateralisation
The cerebrum has a contralateral organisation with each hemisphere of the brain interacting primarily with one half of the body: the left side of the brain interacts with the right side of the body, and vice versa. This is theorized to be caused by a developmental axial twist. Motor connections from the brain to the spinal cord, and sensory connections from the spinal cord to the brain, both cross sides in the brainstem. Visual input follows a more complex rule: the optic nerves from the two eyes come together at a point called the optic chiasm, and half of the fibres from each nerve split off to join the other. The result is that connections from the left half of the retina, in both eyes, go to the left side of the brain, whereas connections from the right half of the retina go to the right side of the brain. Because each half of the retina receives light coming from the opposite half of the visual field, the functional consequence is that visual input from the left side of the world goes to the right side of the brain, and vice versa. Thus, the right side of the brain receives somatosensory input from the left side of the body, and visual input from the left side of the visual field.
The left and right sides of the brain appear symmetrical, but they function asymmetrically. For example, the counterpart of the left-hemisphere motor area controlling the right hand is the right-hemisphere area controlling the left hand. There are, however, several important exceptions, involving language and spatial cognition. The left frontal lobe is dominant for language. If a key language area in the left hemisphere is damaged, it can leave the victim unable to speak or understand, whereas equivalent damage to the right hemisphere would cause only minor impairment to language skills.
A substantial part of current understanding of the interactions between the two hemispheres has come from the study of "split-brain patients"—people who underwent surgical transection of the corpus callosum in an attempt to reduce the severity of epileptic seizures. These patients do not show unusual behaviour that is immediately obvious, but in some cases can behave almost like two different people in the same body, with the right hand taking an action and then the left hand undoing it. These patients, when briefly shown a picture on the right side of the point of visual fixation, are able to describe it verbally, but when the picture is shown on the left, are unable to describe it, but may be able to give an indication with the left hand of the nature of the object shown.
Emotion
Emotions are generally defined as two-step multicomponent processes involving elicitation, followed by psychological feelings, appraisal, expression, autonomic responses, and action tendencies. Attempts to localise basic emotions to certain brain regions have been controversial; some research found no evidence for specific locations corresponding to emotions, but instead found circuitry involved in general emotional processes. The amygdala, orbitofrontal cortex, mid and anterior insular cortex and lateral prefrontal cortex, appeared to be involved in generating the emotions, while weaker evidence was found for the ventral tegmental area, ventral pallidum and nucleus accumbens in incentive salience. Others, however, have found evidence of activation of specific regions, such as the basal ganglia in happiness, the subcallosal cingulate cortex in sadness, and amygdala in fear.
Cognition
The brain is responsible for cognition, which functions through numerous processes and executive functions. Executive functions include the ability to filter information and tune out irrelevant stimuli with attentional control and cognitive inhibition, the ability to process and manipulate information held in working memory, the ability to think about multiple concepts simultaneously and switch tasks with cognitive flexibility, the ability to inhibit impulses and prepotent responses with inhibitory control, and the ability to determine the relevance of information or appropriateness of an action. Higher order executive functions require the simultaneous use of multiple basic executive functions, and include planning, prospection and fluid intelligence (i.e., reasoning and problem solving).
The prefrontal cortex plays a significant role in mediating executive functions. Planning involves activation of the dorsolateral prefrontal cortex (DLPFC), anterior cingulate cortex, angular prefrontal cortex, right prefrontal cortex, and supramarginal gyrus. Working memory manipulation involves the DLPFC, inferior frontal gyrus, and areas of the parietal cortex. Inhibitory control involves multiple areas of the prefrontal cortex, as well as the caudate nucleus and subthalamic nucleus.
Physiology
Neurotransmission
Brain activity is made possible by the interconnections of neurons that are linked together to reach their targets. A neuron consists of a cell body, axon, and dendrites. Dendrites are often extensive branches that receive information in the form of signals from the axon terminals of other neurons. The signals received may cause the neuron to initiate an action potential (an electrochemical signal or nerve impulse) which is sent along its axon to the axon terminal, to connect with the dendrites or with the cell body of another neuron. An action potential is initiated at the initial segment of an axon, which contains a specialised complex of proteins. When an action potential reaches the axon terminal it triggers the release of a neurotransmitter at a synapse that propagates a signal that acts on the target cell. These chemical neurotransmitters include dopamine, serotonin, GABA, glutamate, and acetylcholine. GABA is the major inhibitory neurotransmitter in the brain, and glutamate is the major excitatory neurotransmitter. Neurons link at synapses to form neural pathways, neural circuits, and large elaborate network systems such as the salience network and the default mode network, and the activity between them is driven by the process of neurotransmission.
Metabolism
The brain consumes up to 20% of the energy used by the human body, more than any other organ. In humans, blood glucose is the primary source of energy for most cells and is critical for normal function in a number of tissues, including the brain. The human brain consumes approximately 60% of blood glucose in fasted, sedentary individuals. Brain metabolism normally relies upon blood glucose as an energy source, but during times of low glucose (such as fasting, endurance exercise, or limited carbohydrate intake), the brain uses ketone bodies for fuel with a smaller need for glucose. The brain can also utilize lactate during exercise. The brain stores glucose in the form of glycogen, albeit in significantly smaller amounts than that found in the liver or skeletal muscle. Long-chain fatty acids cannot cross the blood–brain barrier, but the liver can break these down to produce ketone bodies. However, short-chain fatty acids (e.g., butyric acid, propionic acid, and acetic acid) and the medium-chain fatty acids, octanoic acid and heptanoic acid, can cross the blood–brain barrier and be metabolised by brain cells.
Although the human brain represents only 2% of the body weight, it receives 15% of the cardiac output, 20% of total body oxygen consumption, and 25% of total body glucose utilization. The brain mostly uses glucose for energy, and deprivation of glucose, as can happen in hypoglycemia, can result in loss of consciousness. The energy consumption of the brain does not vary greatly over time, but active regions of the cortex consume somewhat more energy than inactive regions, which forms the basis for the functional neuroimaging methods of PET and fMRI. These techniques provide a three-dimensional image of metabolic activity. A preliminary study showed that brain metabolic requirements in humans peak at about five years old.
The function of sleep is not fully understood; however, there is evidence that sleep enhances the clearance of metabolic waste products, some of which are potentially neurotoxic, from the brain and may also permit repair. Evidence suggests that the increased clearance of metabolic waste during sleep occurs via increased functioning of the glymphatic system. Sleep may also have an effect on cognitive function by weakening unnecessary connections.
Research
The brain is not fully understood, and research is ongoing. Neuroscientists, along with researchers from allied disciplines, study how the human brain works. The boundaries between the specialties of neuroscience, neurology and other disciplines such as psychiatry have faded as they are all influenced by basic research in neuroscience.
Neuroscience research has expanded considerably. The "Decade of the Brain", an initiative of the United States Government in the 1990s, is considered to have marked much of this increase in research, and was followed in 2013 by the BRAIN Initiative. The Human Connectome Project was a five-year study launched in 2009 to analyse the anatomical and functional connections of parts of the brain, and has provided much data.
An emerging phase in research may be that of simulating brain activity.
Methods
Information about the structure and function of the human brain comes from a variety of experimental methods, including animals and humans. Information about brain trauma and stroke has provided information about the function of parts of the brain and the effects of brain damage. Neuroimaging is used to visualise the brain and record brain activity. Electrophysiology is used to measure, record and monitor the electrical activity of the cortex. Measurements may be of local field potentials of cortical areas, or of the activity of a single neuron. An electroencephalogram can record the electrical activity of the cortex using electrodes placed non-invasively on the scalp.
Invasive measures include electrocorticography, which uses electrodes placed directly on the exposed surface of the brain. This method is used in cortical stimulation mapping, used in the study of the relationship between cortical areas and their systemic function. By using much smaller microelectrodes, single-unit recordings can be made from a single neuron that give a high spatial resolution and high temporal resolution. This has enabled the linking of brain activity to behaviour, and the creation of neuronal maps.
The development of cerebral organoids has opened ways for studying the growth of the brain, and of the cortex, and for understanding disease development, offering further implications for therapeutic applications.
Imaging
Functional neuroimaging techniques show changes in brain activity that relate to the function of specific brain areas. One technique is functional magnetic resonance imaging (fMRI) which has the advantages over earlier methods of SPECT and PET of not needing the use of radioactive materials and of offering a higher resolution. Another technique is functional near-infrared spectroscopy. These methods rely on the haemodynamic response that shows changes in brain activity in relation to changes in blood flow, useful in mapping functions to brain areas. Resting state fMRI
looks at the interaction of brain regions whilst the brain is not performing a specific task. This is also used to show the default mode network.
Any electrical current generates a magnetic field; neural oscillations induce weak magnetic fields, and in functional magnetoencephalography the current produced can show localised brain function in high resolution. Tractography uses MRI and image analysis to create 3D images of the nerve tracts of the brain. Connectograms give a graphical representation of the neural connections of the brain.
Differences in brain structure can be measured in some disorders, notably schizophrenia and dementia. Different biological approaches using imaging have given more insight for example into the disorders of depression and obsessive-compulsive disorder. A key source of information about the function of brain regions is the effects of damage to them.
Advances in neuroimaging have enabled objective insights into mental disorders, leading to faster diagnosis, more accurate prognosis, and better monitoring.
Gene and protein expression
Bioinformatics is a field of study that includes the creation and advancement of databases, and computational and statistical techniques, that can be used in studies of the human brain, particularly in the areas of gene and protein expression. Bioinformatics and studies in genomics, and functional genomics, generated the need for DNA annotation, a transcriptome technology, identifying genes, their locations and functions. GeneCards is a major database.
, just under 20,000 protein-coding genes are seen to be expressed in the human, and some 400 of these genes are brain-specific. The data that has been provided on gene expression in the brain has fuelled further research into a number of disorders. The long term use of alcohol for example, has shown altered gene expression in the brain, and cell-type specific changes that may relate to alcohol use disorder. These changes have been noted in the synaptic transcriptome in the prefrontal cortex, and are seen as a factor causing the drive to alcohol dependence, and also to other substance abuses.
Other related studies have also shown evidence of synaptic alterations and their loss, in the ageing brain. Changes in gene expression alter the levels of proteins in various neural pathways and this has been shown to be evident in synaptic contact dysfunction or loss. This dysfunction has been seen to affect many structures of the brain and has a marked effect on inhibitory neurons resulting in a decreased level of neurotransmission, and subsequent cognitive decline and disease.
Clinical significance
Injury
Injury to the brain can manifest in many ways. Traumatic brain injury, for example received in contact sport, after a fall, or a traffic or work accident, can be associated with both immediate and longer-term problems. Immediate problems may include bleeding within the brain, this may compress the brain tissue or damage its blood supply. Bruising to the brain may occur. Bruising may cause widespread damage to the nerve tracts that can lead to a condition of diffuse axonal injury. A fractured skull, injury to a particular area, deafness, and concussion are also possible immediate developments. In addition to the site of injury, the opposite side of the brain may be affected, termed a contrecoup injury. Longer-term issues that may develop include posttraumatic stress disorder, and hydrocephalus. Chronic traumatic encephalopathy can develop following multiple head injuries.
Disease
Neurodegenerative diseases result in progressive damage to, or loss of neurons affecting different functions of the brain, that worsen with age. Common types are dementias including Alzheimer's disease, alcoholic dementia, vascular dementia, and Parkinson's disease dementia. Other rarer infectious, genetic, or metabolic types include Huntington's disease, motor neuron diseases, HIV dementia, syphilis-related dementia and Wilson's disease. Neurodegenerative diseases can affect different parts of the brain, and can affect movement, memory, and cognition. Rare prion diseases including Creutzfeldt–Jakob disease and its variant, and kuru are fatal neurodegenerative diseases.
Cerebral atherosclerosis is atherosclerosis that affects the brain. It results from the build-up of plaques formed of cholesterol, in the large arteries of the brain, and can be mild to significant. When significant, arteries can become narrowed enough to reduce blood flow. It contributes to the development of dementia, and has protein similarities to those found in Alzheimer's disease.
The brain, although protected by the blood–brain barrier, can be affected by infections including viruses, bacteria and fungi. Infection may be of the meninges (meningitis), the brain matter (encephalitis), or within the brain matter (such as a cerebral abscess).
Tumours
Brain tumours can be either benign or cancerous. Most malignant tumours arise from another part of the body, most commonly from the lung, breast and skin. Cancers of brain tissue can also occur, and originate from any tissue in and around the brain. Meningioma, cancer of the meninges around the brain, is more common than cancers of brain tissue. Cancers within the brain may cause symptoms related to their size or position, with symptoms including headache and nausea, or the gradual development of focal symptoms such as gradual difficulty seeing, swallowing, talking, or as a change of mood. Cancers are in general investigated through the use of CT scans and MRI scans. A variety of other tests including blood tests and lumbar puncture may be used to investigate for the cause of the cancer and evaluate the type and stage of the cancer. The corticosteroid dexamethasone is often given to decrease the swelling of brain tissue around a tumour. Surgery may be considered, however given the complex nature of many tumours or based on tumour stage or type, radiotherapy or chemotherapy may be considered more suitable.
Mental disorders
Mental disorders, such as depression, schizophrenia, bipolar disorder, posttraumatic stress disorder, attention deficit hyperactivity disorder, obsessive-compulsive disorder, Tourette syndrome, and addiction, are known to relate to the functioning of the brain. Treatment for mental disorders may include psychotherapy, psychiatry, social intervention and personal recovery work or cognitive behavioural therapy; the underlying issues and associated prognoses vary significantly between individuals.
Epilepsy
Epileptic seizures are thought to relate to abnormal electrical activity. Seizure activity can manifest as absence of consciousness, focal effects such as limb movement or impediments of speech, or be generalized in nature. Status epilepticus refers to a seizure or series of seizures that have not terminated within five minutes. Seizures have a large number of causes, however many seizures occur without a definitive cause being found. In a person with epilepsy, risk factors for further seizures may include sleeplessness, drug and alcohol intake, and stress. Seizures may be assessed using blood tests, EEG and various medical imaging techniques based on the medical history and medical examination findings. In addition to treating an underlying cause and reducing exposure to risk factors, anticonvulsant medications can play a role in preventing further seizures.
Congenital
Some brain disorders, such as Tay–Sachs disease, are congenital and linked to genetic and chromosomal mutations. A rare group of congenital cephalic disorders known as lissencephaly is characterised by the lack of, or inadequacy of, cortical folding. Normal development of the brain can be affected during pregnancy by nutritional deficiencies, teratogens, infectious diseases, and by the use of recreational drugs, including alcohol (which may result in fetal alcohol spectrum disorders).
Most cerebral arteriovenous malformations are congenital, these tangled networks of blood vessels may remain without symptoms but at their worst may rupture and cause intracranial hemorrhaging.
Stroke
A stroke is a decrease in blood supply to an area of the brain causing cell death and brain injury. This can lead to a wide range of symptoms, including the "FAST" symptoms of facial droop, arm weakness, and speech difficulties (including with speaking and finding words or forming sentences). Symptoms relate to the function of the affected area of the brain and can point to the likely site and cause of the stroke. Difficulties with movement, speech, or sight usually relate to the cerebrum, whereas imbalance, double vision, vertigo and symptoms affecting more than one side of the body usually relate to the brainstem or cerebellum.
Most strokes result from loss of blood supply, typically because of an embolus, rupture of a fatty plaque causing thrombus, or narrowing of small arteries. Strokes can also result from bleeding within the brain. Transient ischaemic attacks (TIAs) are strokes in which symptoms resolve within 24 hours. Investigation into the stroke will involve a medical examination (including a neurological examination) and the taking of a medical history, focusing on the duration of the symptoms and risk factors (including high blood pressure, atrial fibrillation, and smoking). Further investigation is needed in younger patients. An ECG and biotelemetry may be conducted to identify atrial fibrillation; an ultrasound can investigate narrowing of the carotid arteries; an echocardiogram can be used to look for clots within the heart, diseases of the heart valves or the presence of a patent foramen ovale. Blood tests are routinely done as part of the workup including diabetes tests and a lipid profile.
Some treatments for stroke are time-critical. These include clot dissolution or surgical removal of a clot for ischaemic strokes, and decompression for haemorrhagic strokes. As stroke is time critical, hospitals and even pre-hospital care of stroke involves expedited investigations – usually a CT scan to investigate for a haemorrhagic stroke and a CT or MR angiogram to evaluate arteries that supply the brain. MRI scans, not as widely available, may be able to demonstrate the affected area of the brain more accurately, particularly with ischaemic stroke.
Having experienced a stroke, a person may be admitted to a stroke unit, and treatments may be directed as preventing future strokes, including ongoing anticoagulation (such as aspirin or clopidogrel), antihypertensives, and lipid-lowering drugs. A multidisciplinary team including speech pathologists, physiotherapists, occupational therapists, and psychologists plays a large role in supporting a person affected by a stroke and their rehabilitation. A history of stroke increases the risk of developing dementia by around 70%, and recent stroke increases the risk by around 120%.
Brain death
Brain death refers to an irreversible total loss of brain function. This is characterised by coma, loss of reflexes, and apnoea, however, the declaration of brain death varies geographically and is not always accepted. In some countries there is also a defined syndrome of brainstem death. Declaration of brain death can have profound implications as the declaration, under the principle of medical futility, will be associated with the withdrawal of life support, and as those with brain death often have organs suitable for organ donation. The process is often made more difficult by poor communication with patients' families.
When brain death is suspected, reversible differential diagnoses such as, electrolyte, neurological and drug-related cognitive suppression need to be excluded. Testing for reflexes can be of help in the decision, as can the absence of response and breathing. Clinical observations, including a total lack of responsiveness, a known diagnosis, and neural imaging evidence, may all play a role in the decision to pronounce brain death.
Society and culture
Neuroanthropology is the study of the relationship between culture and the brain. It explores how the brain gives rise to culture, and how culture influences brain development. Cultural differences and their relation to brain development and structure are researched in different fields.
The mind
The philosophy of the mind studies such issues as the problem of understanding consciousness and the mind–body problem. The relationship between the brain and the mind is a significant challenge both philosophically and scientifically. This is because of the difficulty in explaining how mental activities, such as thoughts and emotions, can be implemented by physical structures such as neurons and synapses, or by any other type of physical mechanism. This difficulty was expressed by Gottfried Leibniz in the analogy known as Leibniz's Mill:
Doubt about the possibility of a mechanistic explanation of thought drove René Descartes, and most other philosophers along with him, to dualism: the belief that the mind is to some degree independent of the brain. There has always, however, been a strong argument in the opposite direction. There is clear empirical evidence that physical manipulations of, or injuries to, the brain (for example by drugs or by lesions, respectively) can affect the mind in potent and intimate ways. In the 19th century, the case of Phineas Gage, a railway worker who was injured by a stout iron rod passing through his brain, convinced both researchers and the public that cognitive functions were localised in the brain. Following this line of thinking, a large body of empirical evidence for a close relationship between brain activity and mental activity has led most neuroscientists and contemporary philosophers to be materialists, believing that mental phenomena are ultimately the result of, or reducible to, physical phenomena.
Brain size
The size of the brain and a person's intelligence are not strongly related. Studies tend to indicate small to moderate correlations (averaging around 0.3 to 0.4) between brain volume and IQ. The most consistent associations are observed within the frontal, temporal, and parietal lobes, the hippocampi, and the cerebellum, but these only account for a relatively small amount of variance in IQ, which itself has only a partial relationship to general intelligence and real-world performance.
Other animals, including whales and elephants, have larger brains than humans. However, when the brain-to-body mass ratio is taken into account, the human brain is almost twice as large as that of a bottlenose dolphin, and three times as large as that of a chimpanzee. However, a high ratio does not of itself demonstrate intelligence: very small animals have high ratios and the treeshrew has the largest quotient of any mammal.
In popular culture
Earlier ideas about the relative importance of the different organs of the human body sometimes emphasised the heart.
Modern Western popular conceptions, in contrast, have placed increasing focus on the brain.
Research has disproved some common misconceptions about the brain. These include both ancient and modern myths. It is not true (for example) that neurons are not replaced after the age of two; nor that normal humans use only ten per cent of the brain. Popular culture has also oversimplified the lateralisation of the brain by suggesting that functions are completely specific to one side of the brain or the other. Akio Mori coined the term "game brain" for the unreliably supported theory that spending long periods playing video games harmed the brain's pre-frontal region, and impaired the expression of emotion and creativity.
Historically, particularly in the early-19th century, the brain featured in popular culture through phrenology, a pseudoscience that assigned personality attributes to different regions of the cortex. The cortex remains important in popular culture as covered in books and satire.
The human brain can feature in science fiction, with themes such as brain transplants and cyborgs (beings with features like partly artificial brains). The 1942 science-fiction book (adapted three times for the cinema) Donovan's Brain tells the tale of an isolated brain kept alive in vitro, gradually taking over the personality of the book's protagonist.
History
Early history
The Edwin Smith Papyrus, an ancient Egyptian medical treatise written in the 17th century BC, contains the earliest recorded reference to the brain. The hieroglyph for brain, occurring eight times in this papyrus, describes the symptoms, diagnosis, and prognosis of two traumatic injuries to the head. The papyrus mentions the external surface of the brain, the effects of injury (including seizures and aphasia), the meninges, and cerebrospinal fluid.
In the fifth century BC, Alcmaeon of Croton in Magna Grecia, first considered the brain to be the seat of the mind. Also in the fifth century BC in Athens, the unknown author of On the Sacred Disease, a medical treatise which is part of the Hippocratic Corpus and traditionally attributed to Hippocrates, believed the brain to be the seat of intelligence. Aristotle, in his biology initially believed the heart to be the seat of intelligence, and saw the brain as a cooling mechanism for the blood. He reasoned that humans are more rational than the beasts because, among other reasons, they have a larger brain to cool their hot-bloodedness. Aristotle did describe the meninges and distinguished between the cerebrum and cerebellum.
Herophilus of Chalcedon in the fourth and third centuries BC distinguished the cerebrum and the cerebellum, and provided the first clear description of the ventricles; and with Erasistratus of Ceos experimented on living brains. Their works are now mostly lost, and we know about their achievements due mostly to secondary sources. Some of their discoveries had to be re-discovered a millennium after their deaths. Anatomist physician Galen in the second century AD, during the time of the Roman Empire, dissected the brains of sheep, monkeys, dogs, and pigs. He concluded that, as the cerebellum was denser than the brain, it must control the muscles, while as the cerebrum was soft, it must be where the senses were processed. Galen further theorised that the brain functioned by movement of animal spirits through the ventricles.
Renaissance
In 1316, Mondino de Luzzi's Anathomia began the modern study of brain anatomy.
Niccolò Massa discovered in 1536 that the ventricles were filled with fluid. Archangelo Piccolomini of Rome was the first to distinguish between the cerebrum and cerebral cortex. In 1543 Andreas Vesalius published his seven-volume De humani corporis fabrica. The seventh book covered the brain and eye, with detailed images of the ventricles, cranial nerves, pituitary gland, meninges, structures of the eye, the vascular supply to the brain and spinal cord, and an image of the peripheral nerves. Vesalius rejected the common belief that the ventricles were responsible for brain function, arguing that many animals have a similar ventricular system to humans, but no true intelligence.
René Descartes proposed the theory of dualism to tackle the issue of the brain's relation to the mind. He suggested that the pineal gland was where the mind interacted with the body, serving as the seat of the soul and as the connection through which animal spirits passed from the blood into the brain. This dualism likely provided impetus for later anatomists to further explore the relationship between the anatomical and functional aspects of brain anatomy.
Thomas Willis is considered a second pioneer in the study of neurology and brain science. He wrote Cerebri Anatome () in 1664, followed by Cerebral Pathology in 1667. In these he described the structure of the cerebellum, the ventricles, the cerebral hemispheres, the brainstem, and the cranial nerves, studied its blood supply; and proposed functions associated with different areas of the brain. The circle of Willis was named after his investigations into the blood supply of the brain, and he was the first to use the word "neurology". Willis removed the brain from the body when examining it, and rejected the commonly held view that the cortex only consisted of blood vessels, and the view of the last two millennia that the cortex was only incidentally important.
In the middle of 19th century Emil du Bois-Reymond and Hermann von Helmholtz were able to use a galvanometer to show that electrical impulses passed at measurable speeds along nerves, refuting the view of their teacher Johannes Peter Müller that the nerve impulse was a vital function that could not be measured. Richard Caton in 1875 demonstrated electrical impulses in the cerebral hemispheres of rabbits and monkeys. In the 1820s, Jean Pierre Flourens pioneered the experimental method of damaging specific parts of animal brains describing the effects on movement and behavior.
Modern period
Studies of the brain became more sophisticated with the use of the microscope and the development of a silver staining method by Camillo Golgi during the 1880s. This was able to show the intricate structures of single neurons. This was used by Santiago Ramón y Cajal and led to the formation of the neuron doctrine, the then revolutionary hypothesis that the neuron is the functional unit of the brain. He used microscopy to uncover many cell types, and proposed functions for the cells he saw. For this, Golgi and Cajal are considered the founders of twentieth century neuroscience, both sharing the Nobel prize in 1906 for their studies and discoveries in this field.
Charles Sherrington published his influential 1906 work The Integrative Action of the Nervous System examining the function of reflexes, evolutionary development of the nervous system, functional specialisation of the brain, and layout and cellular function of the central nervous system. In 1942 he coined the term enchanted loom as a metaphor for the brain. John Farquhar Fulton, founded the Journal of Neurophysiology and published the first comprehensive textbook on the physiology of the nervous system during 1938. Neuroscience during the twentieth century began to be recognised as a distinct unified academic discipline, with David Rioch, Francis O. Schmitt, and Stephen Kuffler playing critical roles in establishing the field. Rioch originated the integration of basic anatomical and physiological research with clinical psychiatry at the Walter Reed Army Institute of Research, starting in the 1950s. During the same period, Schmitt established the Neuroscience Research Program, an inter-university and international organisation, bringing together biology, medicine, psychological and behavioural sciences. The word neuroscience itself arises from this program.
Paul Broca associated regions of the brain with specific functions, in particular language in Broca's area, following work on brain-damaged patients. John Hughlings Jackson described the function of the motor cortex by watching the progression of epileptic seizures through the body. Carl Wernicke described a region associated with language comprehension and production. Korbinian Brodmann divided regions of the brain based on the appearance of cells. By 1950, Sherrington, Papez, and MacLean had identified many of the brainstem and limbic system functions. The capacity of the brain to re-organise and change with age, and a recognised critical development period, were attributed to neuroplasticity, pioneered by Margaret Kennard, who experimented on monkeys during the 1930-40s.
Harvey Cushing (1869–1939) is recognised as the first proficient brain surgeon in the world. In 1937, Walter Dandy began the practice of vascular neurosurgery by performing the first surgical clipping of an intracranial aneurysm.
Comparative anatomy
The human brain has many properties that are common to all vertebrate brains. Many of its features are common to all mammalian brains, most notably a six-layered cerebral cortex and a set of associated structures, including the hippocampus and amygdala. The cortex is proportionally larger in humans than in many other mammals. Humans have more association cortex, sensory and motor parts than smaller mammals such as the rat and the cat.
As a primate brain, the human brain has a much larger cerebral cortex, in proportion to body size, than most mammals, and a highly developed visual system.
As a hominid brain, the human brain is substantially enlarged even in comparison to the brain of a typical monkey. The sequence of human evolution from Australopithecus (four million years ago) to Homo sapiens (modern humans) was marked by a steady increase in brain size. As brain size increased, this altered the size and shape of the skull, from about 600 cm3 in Homo habilis to an average of about 1520 cm3 in Homo neanderthalensis. Differences in DNA, gene expression, and gene–environment interactions help explain the differences between the function of the human brain and other primates.
| Biology and health sciences | Human anatomy | null |
491164 | https://en.wikipedia.org/wiki/Lima%20bean | Lima bean | A lima bean (Phaseolus lunatus), also commonly known as butter bean, sieva bean, double bean or Madagascar bean, is a legume grown for its edible seeds or beans.
Origin and uses
Phaseolus lunatus is found in Meso- and South America. Two gene pools of cultivated lima beans point to independent domestication events. The Mesoamerican lima bean is distributed in neotropical lowlands, while the other is found in the western Andes. They were discovered in Peru and may have been the first plant that was brought up under civilization by the native farmers.
The Andean domestication took place around 2000 BC and produced a large-seeded variety (lima type), while the second, taking place in Mesoamerica around 800 AD, produced a small-seeded variety (Sieva type). By around 1300, cultivation had spread north of the Rio Grande, and, in the 1500s, the plant began to be cultivated in the Old World.
The small-seeded (Sieva) type is found distributed from Mexico to Argentina, generally below above sea level, while the large-seeded wild form (lima type) is found distributed in the north of Peru, from above sea level.
The Moche culture (1–800 CE) cultivated lima beans heavily and often depicted them in their art. During the Spanish Viceroyalty of Peru, lima beans were exported to the rest of the Americas and Europe, and since the boxes of such goods had their place of origin labeled "Lima, Peru", the beans got named as such.
The term "butter bean" is widely used in North and South Carolina for a large, flat and yellow/white variety of lima bean (P. lunatus var. macrocarpus, or P. limensis). In the United States, Sieva-type beans are traditionally called butter beans, also otherwise known as the Dixie or Henderson type. In that area, lima beans and butter beans are seen as two distinct types of beans, although they are the same species. In the United Kingdom and the United States, "butter beans" refers to either dried beans, which can be purchased to rehydrate, or the canned variety, which are ready to use. In culinary use there, lima beans and butter beans are distinct, the former being small and green, the latter large and yellow. In areas where both are considered to be lima beans, the green variety may be labeled as "baby" (and less commonly "junior") limas.
In Spain, it is called garrofón and constitutes one of the main ingredients of the famous Valencian paella.
In India, they are called double beans. Dried beans are soaked overnight and pressure-cooked as ingredients in curries.
Domestication
The lima bean is a domesticated species of economic and cultural importance worldwide, especially in Mexico. The species has two varieties. The wild variety is silvester and the domesticated one is lunatus.
Crop
In the U.S., it is a warm-season crop, grown mainly in Delaware and the mid-Atlantic region for processing and in the Midwest and California for dry beans. Baby lima beans are planted in early June and harvested about 10–12 weeks later. In western New York State, baby lima bean production increased greatly from 2011 to 2015.
Cultivation and cultivars
Cultivation
In Oaxaca, Mexico, the main rainy season lasts from June to August, and most of the above-ground parts die during the dry season. Germination or budding occurs in June or July. The first inflorescence is in October or November. The production of flowers and fruits usually ends between February and April.
Cultivars
Both bush and pole (vine) cultivars exist; the latter range from in height. The bush cultivars mature earlier than the pole cultivars. The pods are up to long. The mature seeds are long and oval to kidney-shaped. In most cultivars, the seeds are quite flat, but in the "potato" cultivars, the shape approaches spherical. White seeds are common, but black, red, orange, and variously mottled seeds are also known. The immature seeds are uniformly green. Lima beans typically yield of seed and of biomass per hectare.
The seeds of the cultivars listed below are white unless otherwise noted. Closely related or synonymous names are listed on the same line.
Bush types
'Henderson' / 'Thorogreen', 65 days (heirloom)
'Eastland', 68 days
'Jackson Wonder', 68 days (heirloom, seeds brown mottled with purple)
'Dixie Butterpea', 75 days (heirloom, two strains are common: red speckled and white seeded)
'Fordhook 242', 75 days, 1945 AAS winner
Pole types
'Carolina Sieva', 75 days (heirloom, suffered a seed crop failure in the years 2011 and 2012, causing this variety to still not be widely sourced ten years later)
'Christmas' / 'Chestnut' / 'Giant Speckled' / 'Speckled Calico', 78 days (heirloom, seeds white mottled with red)
'Big 6' / 'Big Mama', 80 days
'Willow Leaf', 80 days (heirloom, there are white-seeded and variously mottled strains)
'Mezcla', 82 days
'King of the Garden', 85 days (heirloom)
Pathogens/disease
Phytophthora phaseoli is one example of a pathogen of the lima bean. It is an oomycete plant pathogen that causes downy mildew of lima bean during cool and humid weather conditions. To combat this pathogen, developing lima bean cultivars with resistance is a relatively cost-efficient method that is also environmentally safe as compared to using pesticides.
Didymella is a foliar disease found in baby lima beans first reported in New York State. Symptoms include small necrotic tan spots with red to reddish brown irregular margins that come together to eventually cover the entire leaf. Lesions occur after around 3–4 weeks of planting and increase until there is considerable defoliation. Lesions are usually observed on the stems. Two pycnidial fungi were found on leaves, including Didymella sp. and Boeremia exigua var. exigua, which is pathogenic on baby lima bean and plays a role in the foliar disease complex. Other fungal diseases on lima beans with similar symptoms are B. exigua var. exigua, pod blight caused by Diaporthe phaseolorum, and leaf spots caused by Phyllosticta sp. and Phoma subcircinata.
Predators/hosts
The two-spotted spider mites or Tetranychus urticae lay eggs on lima bean leaves. It prefers lima bean plants as a host food source over other plants such as tomato or cabbage plants.
Spider mites pose the greatest threat to lima bean plants compared to other species, such as the Common cutworm (Spodoptera litura), which is also known to feed on lima bean plants. These plants are host plants for their larvae.
One herbivore of lima bean is Spodoptera littoralis, the African cotton leafworm. An attack by this herbivore induces hydrogen peroxide in the leaves. This may also be advantageous to defend against pathogens such as bacteria, fungi, or viruses, as they can easily invade herbivore-infected leaves.
Other predatory insects include ants, wasps, flies and beetles.
Defenses
Lima beans use extrafloral nectar (EFN) secretion when exposed to volatiles from other plants infested by herbivore species. Producing EFN can be an indirect defense since it supplies enemies of herbivores with an alternative food source. The predator of lima beans, spider mites, also have their own predators, the carnivorous mite Phytoseiulus persimilis. These predatory mites use EFN as an alternative food source, and thus, the production of this by the lima bean can attract P. persimilis and thus deter their herbivore hosts.
The main induced defense of the lima bean is the Jasmonic acid pathway. Jasmonic acid induces the production of extrafloral nectar flow or induces it when herbivory occurs, such as when attacked by spider mites.
One direct chemical defense involves cyanogenesis, which is the release of hydrogen cyanide when the cell senses damage. Cyanide acts as a repellent on the leaves of lima beans.
Plant behavior
Phaseolus lunatus has adapted to live in many different climates around the world. One of these adaptations includes a particularly effective induced herbivory defense. The lima bean is able to signal to the carnivorous natural enemy of herbivores, the carnivorous mite, mediated by HIPVs (Herbivore Induced Plant Volatiles) in an attempt to save itself from further predation.
The lima bean plant does this as an induced defense when being eaten by herbivorous predators. It is the mechanical wounding and chemical elicitors from insect oral secretions that first begin the signaling pathway to induce HIPV production. Once this pathway is induced, the plant produces HIPVs which are released into the air and can be received by any organisms that have receptors capable of receiving HIPVs, which includes: carnivores, conspecific and heterospecific herbivores, as well as neighboring plants. It is this signaling of the carnivorous natural enemy of herbivores that is of particular interest, as they become attracted to the plant and will then come and prey upon the plant's herbivorous enemy, thereby reducing herbivory of the plant.
One particular experiment in which this was made apparent was in the understanding of the tritrophic system between the lima bean plant, two-spotted spider mite, and the carnivorous mite. Here, experimenters noticed an increase in HIPVs when the lima bean plant was preyed on by the two-spotted spider mite. Then, when the carnivorous mite was introduced, it had increased prey-searching efficacy and overall attraction to the lima bean plant, even once the two-spotted spider mite was removed, but the HIPVs remained high.
Toxicity
Like many beans, raw lima beans are toxic (containing e.g. phytohaemagglutinin) if not boiled for at least 10 minutes. Canned beans can be eaten without having to be boiled first, as they are pre-cooked.
The lima bean can contain anti-nutrients like phytic acids, saponin, oxalate, tannin, and trypsin inhibitor. These inhibit the absorption of nutrients in animals and can cause damage to some organs. In addition to boiling, methods of roasting, pressure cooking, soaking, and germination can also reduce the antinutrients significantly.
Nutrition
The most abundant mineral in the raw lima bean is potassium, followed by calcium, phosphorus, magnesium, sodium, and iron. When lima beans germinate, there is increased bioavailability of calcium and phosphorus. Additionally, it is a good source of vitamin B6.
Uses
Culinary
Like many other legumes, lima beans are a good source of dietary fiber and a virtually fat-free source of high-quality protein.
Lima beans contain both soluble fiber, which helps regulate blood sugar levels and lowers cholesterol, and insoluble fiber, which aids in the prevention of constipation, digestive disorders, irritable bowel syndrome, and diverticulitis.
Medical
Blood sugar level
The high fiber content in lima beans prevents blood sugar levels from rising too rapidly after eating them due to the presence of those large amounts of absorption-slowing compounds in the beans, and the high soluble fiber content. Soluble fiber absorbs water in the stomach, forming a gel that slows down the absorption of the bean's carbohydrates. They can, therefore, help balance blood sugar levels while providing steady, slow-burning energy, which makes them a good choice for people with diabetes suffering from insulin resistance.
| Biology and health sciences | Fabales | null |
8334453 | https://en.wikipedia.org/wiki/Stadion%20%28unit%29 | Stadion (unit) | The stadion (plural stadia, ; latinized as stadium), also anglicized as stade, was an ancient Greek unit of length, consisting of 600 Ancient Greek feet (podes). Its exact length is unknown today; historians estimate it at between 150 m and 210 m.
Calculations
According to Herodotus, one stadium was equal to 600 Greek feet (podes). However, the length of the foot varied in different parts of the Greek world, and the length of the stadion has been the subject of argument and hypothesis for hundreds of years.
An empirical determination of the length of the stadion was made by Lev Vasilevich Firsov, who compared 81 distances given by Eratosthenes and Strabo with the straight-line distances measured by modern methods, and averaged the results. He obtained a result of about . Various equivalent lengths have been proposed, and some have been named. Among them are:
Which measure of the stadion is used can affect the interpretation of ancient texts. For example, the error in the calculation of Earth's circumference by Eratosthenes or Posidonius is dependent on which stadion is chosen to be appropriate.
Other uses
From the Middle Ages on, the word stadium has been used as a synonym for the furlong (which is 220 yards, equal to one eighth of a mile), which is of Old English origin.
| Physical sciences | Other | Basics and measurement |
8335148 | https://en.wikipedia.org/wiki/Unbarred%20spiral%20galaxy | Unbarred spiral galaxy | An unbarred spiral galaxy is a type of spiral galaxy without a central bar, or one that is not a barred spiral galaxy. It is designated with an SA in the galaxy morphological classification scheme. A nearby example of an unbarred spiral is the Triangulum Galaxy.
Barless spiral galaxies are one of three general types of spiral galaxies under the de Vaucouleurs system classification system, the other two being intermediate spiral galaxy and barred spiral galaxy. Under the Hubble tuning fork, it is one of two general types of spiral galaxy, the other being barred spirals.
Grades
Unbarred lenticular galaxy
An unbarred lenticular galaxy is a lenticular version of an unbarred spiral galaxy. They have the Hubble type of SA0.
An example of this is the galaxy AM 0644-741. For other examples, see :Category:Unbarred lenticular galaxies.
| Physical sciences | Galaxy classification | Astronomy |
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