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**Secure Shell**
Secure Shell:
The Secure Shell Protocol (SSH) is a cryptographic network protocol for operating network services securely over an unsecured network. Its most notable applications are remote login and command-line execution.
Secure Shell:
SSH applications are based on a client–server architecture, connecting an SSH client instance with an SSH server. SSH operates as a layered protocol suite comprising three principal hierarchical components: the transport layer provides server authentication, confidentiality, and integrity; the user authentication protocol validates the user to the server; and the connection protocol multiplexes the encrypted tunnel into multiple logical communication channels.SSH was designed on Unix-like operating systems, as a replacement for Telnet and for unsecured remote Unix shell protocols, such as the Berkeley Remote Shell (rsh) and the related rlogin and rexec protocols, which all use insecure, plaintext transmission of authentication tokens.
Secure Shell:
SSH was first designed in 1995 by Finnish computer scientist Tatu Ylönen. Subsequent development of the protocol suite proceeded in several developer groups, producing several variants of implementation. The protocol specification distinguishes two major versions, referred to as SSH-1 and SSH-2. The most commonly implemented software stack is OpenSSH, released in 1999 as open-source software by the OpenBSD developers. Implementations are distributed for all types of operating systems in common use, including embedded systems.
Definition:
SSH uses public-key cryptography to authenticate the remote computer and allow it to authenticate the user, if necessary.SSH may be used in several methodologies. In the simplest manner, both ends of a communication channel use automatically generated public-private key pairs to encrypt a network connection, and then use a password to authenticate the user.
Definition:
When the public-private key pair is generated by the user manually, the authentication is essentially performed when the key pair is created, and a session may then be opened automatically without a password prompt. In this scenario, the public key is placed on all computers that must allow access to the owner of the matching private key, which the owner keeps private. While authentication is based on the private key, the key is never transferred through the network during authentication. SSH only verifies that the same person offering the public key also owns the matching private key.
Definition:
In all versions of SSH it is important to verify unknown public keys, i.e. associate the public keys with identities, before accepting them as valid. Accepting an attacker's public key without validation will authorize an unauthorized attacker as a valid user.
Authentication: OpenSSH key management:
On Unix-like systems, the list of authorized public keys is typically stored in the home directory of the user that is allowed to log in remotely, in the file ~/.ssh/authorized_keys. This file is respected by SSH only if it is not writable by anything apart from the owner and root. When the public key is present on the remote end and the matching private key is present on the local end, typing in the password is no longer required. However, for additional security the private key itself can be locked with a passphrase.
Authentication: OpenSSH key management:
The private key can also be looked for in standard places, and its full path can be specified as a command line setting (the option -i for ssh). The ssh-keygen utility produces the public and private keys, always in pairs.
Authentication: OpenSSH key management:
SSH also supports password-based authentication that is encrypted by automatically generated keys. In this case, the attacker could imitate the legitimate server side, ask for the password, and obtain it (man-in-the-middle attack). However, this is possible only if the two sides have never authenticated before, as SSH remembers the key that the server side previously used. The SSH client raises a warning before accepting the key of a new, previously unknown server. Password authentication can be disabled from the server side.
Use:
SSH is typically used to log into a remote machine and execute commands, but it also supports tunneling, forwarding TCP ports and X11 connections; it can transfer files using the associated SSH file transfer (SFTP) or secure copy (SCP) protocols. SSH uses the client–server model.
Use:
An SSH client program is typically used for establishing connections to an SSH daemon, such as sshd, accepting remote connections. Both are commonly present on most modern operating systems, including macOS, most distributions of Linux, OpenBSD, FreeBSD, NetBSD, Solaris and OpenVMS. Notably, versions of Windows prior to Windows 10 version 1709 do not include SSH by default. Proprietary, freeware and open source (e.g. PuTTY, and the version of OpenSSH which is part of Cygwin) versions of various levels of complexity and completeness exist. File managers for UNIX-like systems (e.g. Konqueror) can use the FISH protocol to provide a split-pane GUI with drag-and-drop. The open source Windows program WinSCP provides similar file management (synchronization, copy, remote delete) capability using PuTTY as a back-end. Both WinSCP and PuTTY are available packaged to run directly off a USB drive, without requiring installation on the client machine. The secure shell extension to the Chrome browser also allows SSH connections without any software installation and even allows SSH from a Chromebook computer. Setting up an SSH server in Windows typically involves enabling a feature in Settings app. In Windows 10 version 1709, an official Win32 port of OpenSSH is available.
Use:
SSH is important in cloud computing to solve connectivity problems, avoiding the security issues of exposing a cloud-based virtual machine directly on the Internet. An SSH tunnel can provide a secure path over the Internet, through a firewall to a virtual machine.The IANA has assigned TCP port 22, UDP port 22 and SCTP port 22 for this protocol. IANA had listed the standard TCP port 22 for SSH servers as one of the well-known ports as early as 2001. SSH can also be run using SCTP rather than TCP as the connection oriented transport layer protocol.
Historical development:
Version 1 In 1995, Tatu Ylönen, a researcher at Helsinki University of Technology, Finland, designed the first version of the protocol (now called SSH-1) prompted by a password-sniffing attack at his university network. The goal of SSH was to replace the earlier rlogin, TELNET, FTP and rsh protocols, which did not provide strong authentication nor guarantee confidentiality. Ylönen released his implementation as freeware in July 1995, and the tool quickly gained in popularity. Towards the end of 1995, the SSH user base had grown to 20000 users in fifty countries.In December 1995, Ylönen founded SSH Communications Security to market and develop SSH. The original version of the SSH software used various pieces of free software, such as GNU libgmp, but later versions released by SSH Communications Security evolved into increasingly proprietary software.
Historical development:
It was estimated that by 2000 the number of users had grown to 2 million.
Historical development:
Version 2 "Secsh" was the official Internet Engineering Task Force's (IETF) name for the IETF working group responsible for version 2 of the SSH protocol. In 2006, a revised version of the protocol, SSH-2, was adopted as a standard. This version is incompatible with SSH-1. SSH-2 features both security and feature improvements over SSH-1. Better security, for example, comes through Diffie–Hellman key exchange and strong integrity checking via message authentication codes. New features of SSH-2 include the ability to run any number of shell sessions over a single SSH connection. Due to SSH-2's superiority and popularity over SSH-1, some implementations such as libssh (v0.8.0+), Lsh and Dropbear support only the SSH-2 protocol.
Historical development:
Version 1.99 In January 2006, well after version 2.1 was established, RFC 4253 specified that an SSH server supporting 2.0 as well as prior versions should identify its protocol version as 1.99. This version number does not reflect a historical software revision, but a method to identify backward compatibility.
Historical development:
OpenSSH and OSSH In 1999, developers, desiring availability of a free software version, restarted software development from the 1.2.12 release of the original SSH program, which was the last released under an open source license. This served as a code base for Björn Grönvall's OSSH software. Shortly thereafter, OpenBSD developers forked Grönvall's code and created OpenSSH, which shipped with Release 2.6 of OpenBSD. From this version, a "portability" branch was formed to port OpenSSH to other operating systems.As of 2005, OpenSSH was the single most popular SSH implementation, being the default version in a large number of operating system distributions. OSSH meanwhile has become obsolete. OpenSSH continues to be maintained and supports the SSH-2 protocol, having expunged SSH-1 support from the codebase in the OpenSSH 7.6 release.
Uses:
SSH is a protocol that can be used for many applications across many platforms including most Unix variants (Linux, the BSDs including Apple's macOS, and Solaris), as well as Microsoft Windows. Some of the applications below may require features that are only available or compatible with specific SSH clients or servers. For example, using the SSH protocol to implement a VPN is possible, but presently only with the OpenSSH server and client implementation.
Uses:
For login to a shell on a remote host (replacing Telnet and rlogin) For executing a single command on a remote host (replacing rsh) For setting up automatic (passwordless) login to a remote server (for example, using OpenSSH) In combination with rsync to back up, copy and mirror files efficiently and securely For forwarding a port For tunneling (not to be confused with a VPN, which routes packets between different networks, or bridges two broadcast domains into one).
Uses:
For using as a full-fledged encrypted VPN. Note that only OpenSSH server and client supports this feature.
For forwarding X from a remote host (possible through multiple intermediate hosts) For browsing the web through an encrypted proxy connection with SSH clients that support the SOCKS protocol.
For securely mounting a directory on a remote server as a filesystem on a local computer using SSHFS.
For automated remote monitoring and management of servers through one or more of the mechanisms discussed above.
For development on a mobile or embedded device that supports SSH.
For securing file transfer protocols.
File transfer protocols The Secure Shell protocols are used in several file transfer mechanisms.
Secure copy (SCP), which evolved from RCP protocol over SSH rsync, intended to be more efficient than SCP. Generally runs over an SSH connection.
SSH File Transfer Protocol (SFTP), a secure alternative to FTP (not to be confused with FTP over SSH or FTPS) Files transferred over shell protocol (FISH), released in 1998, which evolved from Unix shell commands over SSH Fast and Secure Protocol (FASP), aka Aspera, uses SSH for control and UDP ports for data transfer.
Architecture:
The SSH protocol has a layered architecture with three separate components: The transport layer (RFC 4253) typically uses the Transmission Control Protocol (TCP) of TCP/IP, reserving port number 22 as a server listening port. This layer handles initial key exchange as well as server authentication, and sets up encryption, compression, and integrity verification. It exposes to the upper layer an interface for sending and receiving plaintext packets with a size of up to 32,768 bytes each, but more can be allowed by each implementation. The transport layer also arranges for key re-exchange, usually after 1 GB of data has been transferred or after one hour has passed, whichever occurs first.
Architecture:
The user authentication layer (RFC 4252) handles client authentication, and provides a suite of authentication algorithms. Authentication is client-driven: when one is prompted for a password, it may be the SSH client prompting, not the server. The server merely responds to the client's authentication requests. Widely used user-authentication methods include the following: password: a method for straightforward password authentication, including a facility allowing a password to be changed. Not all programs implement this method.
Architecture:
publickey: a method for public-key-based authentication, usually supporting at least DSA, ECDSA or RSA keypairs, with other implementations also supporting X.509 certificates.
Architecture:
keyboard-interactive (RFC 4256): a versatile method where the server sends one or more prompts to enter information and the client displays them and sends back responses keyed-in by the user. Used to provide one-time password authentication such as S/Key or SecurID. Used by some OpenSSH configurations when PAM is the underlying host-authentication provider to effectively provide password authentication, sometimes leading to inability to log in with a client that supports just the plain password authentication method.
Architecture:
GSSAPI authentication methods which provide an extensible scheme to perform SSH authentication using external mechanisms such as Kerberos 5 or NTLM, providing single sign-on capability to SSH sessions. These methods are usually implemented by commercial SSH implementations for use in organizations, though OpenSSH does have a working GSSAPI implementation.
Architecture:
The connection layer (RFC 4254) defines the concept of channels, channel requests, and global requests, which define the SSH services provided. A single SSH connection can be multiplexed into multiple logical channels simultaneously, each transferring data bidirectionally. Channel requests are used to relay out-of-band channel-specific data, such as the changed size of a terminal window, or the exit code of a server-side process. Additionally, each channel performs its own flow control using the receive window size. The SSH client requests a server-side port to be forwarded using a global request. Standard channel types include: shell for terminal shells, SFTP and exec requests (including SCP transfers) direct-tcpip for client-to-server forwarded connections forwarded-tcpip for server-to-client forwarded connections The SSHFP DNS record (RFC 4255) provides the public host key fingerprints in order to aid in verifying the authenticity of the host.This open architecture provides considerable flexibility, allowing the use of SSH for a variety of purposes beyond a secure shell. The functionality of the transport layer alone is comparable to Transport Layer Security (TLS); the user-authentication layer is highly extensible with custom authentication methods; and the connection layer provides the ability to multiplex many secondary sessions into a single SSH connection, a feature comparable to BEEP and not available in TLS.
Algorithms:
EdDSA, ECDSA, RSA and DSA for public-key cryptography.
ECDH and Diffie–Hellman for key exchange.
HMAC, AEAD and UMAC for MAC.
AES (and deprecated RC4, 3DES, DES) for symmetric encryption.
AES-GCM and ChaCha20-Poly1305 for AEAD encryption.
SHA (and deprecated MD5) for key fingerprint.
Vulnerabilities:
SSH-1 In 1998, a vulnerability was described in SSH 1.5 which allowed the unauthorized insertion of content into an encrypted SSH stream due to insufficient data integrity protection from CRC-32 used in this version of the protocol. A fix known as SSH Compensation Attack Detector was introduced into most implementations. Many of these updated implementations contained a new integer overflow vulnerability that allowed attackers to execute arbitrary code with the privileges of the SSH daemon, typically root.
Vulnerabilities:
In January 2001 a vulnerability was discovered that allows attackers to modify the last block of an IDEA-encrypted session. The same month, another vulnerability was discovered that allowed a malicious server to forward a client authentication to another server.Since SSH-1 has inherent design flaws which make it vulnerable, it is now generally considered obsolete and should be avoided by explicitly disabling fallback to SSH-1. Most modern servers and clients support SSH-2.
Vulnerabilities:
CBC plaintext recovery In November 2008, a theoretical vulnerability was discovered for all versions of SSH which allowed recovery of up to 32 bits of plaintext from a block of ciphertext that was encrypted using what was then the standard default encryption mode, CBC. The most straightforward solution is to use CTR, counter mode, instead of CBC mode, since this renders SSH resistant to the attack.
Vulnerabilities:
Suspected decryption by NSA On December 28, 2014 Der Spiegel published classified information leaked by whistleblower Edward Snowden which suggests that the National Security Agency may be able to decrypt some SSH traffic. The technical details associated with such a process were not disclosed. A 2017 analysis of the CIA hacking tools BothanSpy and Gyrfalcon suggested that the SSH protocol was not compromised.
Standards documentation:
The following RFC publications by the IETF "secsh" working group document SSH-2 as a proposed Internet standard.
Standards documentation:
RFC 4250 – The Secure Shell (SSH) Protocol Assigned Numbers RFC 4251 – The Secure Shell (SSH) Protocol Architecture RFC 4252 – The Secure Shell (SSH) Authentication Protocol RFC 4253 – The Secure Shell (SSH) Transport Layer Protocol RFC 4254 – The Secure Shell (SSH) Connection Protocol RFC 4255 – Using DNS to Securely Publish Secure Shell (SSH) Key Fingerprints RFC 4256 – Generic Message Exchange Authentication for the Secure Shell Protocol (SSH) RFC 4335 – The Secure Shell (SSH) Session Channel Break Extension RFC 4344 – The Secure Shell (SSH) Transport Layer Encryption Modes RFC 4345 – Improved Arcfour Modes for the Secure Shell (SSH) Transport Layer ProtocolThe protocol specifications were later updated by the following publications: RFC 4419 – Diffie-Hellman Group Exchange for the Secure Shell (SSH) Transport Layer Protocol (March 2006) RFC 4432 – RSA Key Exchange for the Secure Shell (SSH) Transport Layer Protocol (March 2006) RFC 4462 – Generic Security Service Application Program Interface (GSS-API) Authentication and Key Exchange for the Secure Shell (SSH) Protocol (May 2006) RFC 4716 – The Secure Shell (SSH) Public Key File Format (November 2006) RFC 4819 – Secure Shell Public Key Subsystem (March 2007) RFC 5647 – AES Galois Counter Mode for the Secure Shell Transport Layer Protocol (August 2009) RFC 5656 – Elliptic Curve Algorithm Integration in the Secure Shell Transport Layer (December 2009) RFC 6187 – X.509v3 Certificates for Secure Shell Authentication (March 2011) RFC 6239 – Suite B Cryptographic Suites for Secure Shell (SSH) (May 2011) RFC 6594 – Use of the SHA-256 Algorithm with RSA, Digital Signature Algorithm (DSA), and Elliptic Curve DSA (ECDSA) in SSHFP Resource Records (April 2012) RFC 6668 – SHA-2 Data Integrity Verification for the Secure Shell (SSH) Transport Layer Protocol (July 2012) RFC 7479 – Ed25519 SSHFP Resource Records (March 2015) RFC 5592 – Secure Shell Transport Model for the Simple Network Management Protocol (SNMP) (June 2009) RFC 6242 – Using the NETCONF Protocol over Secure Shell (SSH) (June 2011) RFC 8332 – Use of RSA Keys with SHA-256 and SHA-512 in the Secure Shell (SSH) Protocol (March 2018) draft-gerhards-syslog-transport-ssh-00 – SSH transport mapping for SYSLOG (July 2006) draft-ietf-secsh-filexfer-13 – SSH File Transfer Protocol (July 2006)In addition, the OpenSSH project includes several vendor protocol specifications/extensions: OpenSSH PROTOCOL overview OpenSSH certificate/key overview draft-miller-ssh-agent-04 - SSH Agent Protocol (December 2019) | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Human jaw shrinkage**
Human jaw shrinkage:
Human jaw shrinkage is the phenomenon of continued size reduction of the human mandible and maxilla over the past 12,000 to 15,000 years. Modern human lifestyles and diets are vastly different now from what they were for most of human evolutionary history. Human jaws, as well as oral cavities, have been shrinking ever since the Neolithic agricultural revolution (approximately 12,000 years ago). This has been confirmed by bone remains dated to this time period. Researchers are able to infer the basic lifestyle practices of past cultures, enabling them to link jaw size with lifestyle practice/behaviors. Bones from burial sites of past hunter-gatherer societies are associated with larger jaws and mouths, while bones retrieved from former farming cultures have decreased jaw size. Bones from farming societies also indicate the presence of dental malocclusions, commonly known as non-straight teeth. Within recent centuries, as food has become more processed and soft in form, a rapid increase in non-straight teeth, smaller jaws, and mouths; a lack of space for wisdom teeth; and associated health conditions have been observed. Such conditions include sleep apnea, constricted airways, and decreased respiratory fitness. Medical professionals have been making similar observations and documenting them for hundreds of years. Changes in diet, lifestyle, and breathing patterns have led to maladaptive phenotypic expression in terms of morphological craniofacial development that starts in childhood but persists throughout the lifespan.
Evolutionary history:
The general trend of jaw and oral cavity shrinkage, as well as dental malocclusion presence, has been observed in burial remains across Eurasia. Analyses of remains from areas thought to be in situ (origin) to agriculture, such as those in the Levant region dated to approximately 12,000 years ago, are thought to be where humans first changed from hunting and gathering to a more agricultural lifestyle, with some populations relying on agriculture more than others. Burial sites ranging from 15,000 years ago to approximately 4,000 years ago, spread throughout Europe and modern-day Turkey, have been determined to be the remains of farmers, hunter-gatherers, transitional farmers, and semi-sedentary hunter-gatherers; comparisons and analyses of dental dimensions and jaw morphology have been made between these four lifestyle practices. Clear morphological differences were found based on lifestyle practice, as the jaws and teeth associated with more farming were shown to be smaller on average and often accompanied with malocclusion. Hunter-gatherer populations overwhelmingly had larger jaws, almost always providing adequate space for teeth, including wisdom teeth, and tongue crowding was rare. Indigenous hunter-gatherer populations living today, such as Australian aborigines and the Hadza people of Tanzania, have better oral health and less malocclusion than the average human living in a developed society today. Within Hadza populations, a difference in oral health between males and females has been shown; Hadza women predominantly eat agricultural foods, because they spend most of their time in villages, while the men mostly live in the bush, which consists of hunting, tracking, and gathering. This lifestyle difference causes Hadza men's diets to be dominated by wild foods and leads to them having less periodontal disease, straighter teeth, and fewer cavities than their female counterparts. Even when comparing medieval skulls (approximately 500 to 1,500 years ago) with modern skulls, there is stark contrast in terms of jaw size and malocclusion rates.Due to the exponential increase in advancement since the Agricultural Revolution 12,000 years ago, humans' immediate environments, diets, and culture have changed dramatically. This short length of time, relative to evolutionary timescale, means human genetics are still essentially the same as before these modern changes in lifestyle practices. Specific human developmental pathways were naturally selected for, over vast periods of time; however, these pathways no longer fully match our current environments, leading to the rise in new pathologies and disease; this is also known as evolutionary mismatch.
Etiology:
The main contributing factor to the recent increase in malocclusion is widely considered to be due to a sharp reduction in chewing stress, especially during critical periods of craniofacial growth. Experiments done on non-human subjects have shown that induced nasal blockages and/or dietary changes earlier in life lead to maladaptive morphological change in their jaws, intended to simulate what we are observing globally in human children. Significant craniofacial changes due to diet have even been experimentally shown in pigs during development; researchers fed groups either a hard-consistency diet or a soft-consistency diet, for eight months in total. Drastic differences in jaw and facial musculature, facial structure, and tooth-crowding were observed; researchers directly related the findings to what we are observing more in human populations.
Etiology:
Breathing Mode Orthodontics has also allowed us to identify another contributing factor to shrinking mandibles and overall craniofacial morphological change. An overwhelming proportion of orthodontics patients, who are attempting to correct malocclusion of their teeth, share the characteristic of breathing primarily through the oral cavity. Oftentimes, this habitual mouth breathing is caused by obstructed nasal airways during childhood. Modern humans have spent more time indoors and, as a result, are exposed to higher concentrations of allergens, which accumulate to higher concentrations indoors. Children are experiencing allergies at higher rates, causing congested nasal airways, propagating them to breathe through the mouth more often. Chronic mouth breathing in children has been shown to cause posterior-jaw positioning and more crooked teeth and impacts overall jaw development negatively; These morphological changes further constrict airways, also leaving less room for the tongue to rest, leading to higher rates of obstructive sleep apnea. Decreased mandible size has been directly identified as a risk factor for obstructive sleep apnea. Obstructive sleep apnea in non-obese children has been shown to be a direct result of abnormal oral–facial development, with abnormal development being directly tied to decreased muscle tone of oral and facial muscles. This hypotonia of craniofacial muscles can be caused by lack of chewing stress, jaw posture and rest position, chronic nasal airway obstruction, and even respiratory inefficiency.Nasal breathing has been shown to be advantageous to mouth breathing due to a number of factors, such as how the nasal cavity humidifies incoming air, easing the burden on the lungs, while also filtering out a majority of incoming debris and dust. Nasal breathing also promotes a slower breathing rate. Reduced breathing rates have been shown to promote improved health and longevity. Children who are confirmed clinically to be mouth breathers often show considerably higher rates of concentration difficulties, craniofacial bone abnormalities, malocclusion, cross-bite, chronic gingivitis, candida infections, and halitosis. Due to increasingly sedentary lifestyles, overall population fitness levels are thought to contribute as well. Due to a lack of respiratory efficiency, people are overbreathing through the mouth, even when performing non-strenuous tasks. Breathing chronically through the mouth causes a change in rest posture for the jaw; over time, this can significantly alter jaw development in children, as well as adults to an extent. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Baby bumper headguard cap**
Baby bumper headguard cap:
A baby bumper headguard cap, also known as a falling cap, or pudding hat, is a protective hat worn by children learning to walk, to protect their heads in case of falls.Known as a pudding or black pudding, a version used during the early 17th century until the late 18th century was usually open at the top and featured a sausage-shaped bumper roll that circled the head like a crown. It was fastened with straps under the chin.
Baby bumper headguard cap:
The modern-day version can be many colors and may cover the entire head like a helmet. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Plant LED incubator**
Plant LED incubator:
A plant LED incubator is a chamber which can automatically control the environment of the plant. It can control the temperature, moisture, and especially light regime of the plant based on light emitting diodes (LEDs). LEDs have efficient electric lighting with desired wavelengths (Red+Blue) which support greenhouse production in a minimum time and with high quality and quantity. As LEDs are cool it helps plants to be placed as close as possible to light sources without overheating or scorching. This saves space for intense cultivation. It could provide the opportunity of greenhouse-produced fruits and vegetable to be available for the market more quickly and less expensively due to the effect of LED lighting on earliness, compactness and quality of products.
Plant LED incubator:
Incandescent and fluorescent lamps currently available for lighting greenhouse, phytotrones and plant incubators emit color bands that may cause unwanted stem elongation and low quality in plant species, are electrically inefficient, short-lived, and particularly not eco-friendly for having hazardous waste disposal issues. On the other side, high-pressure sodium (HPS) discharge lamps have been well established in the greenhouses for their sufficient light intensities which support transplants and seedlings growth and development. Unfortunately, they have also many drawbacks. They are intensely hot and scorch nearby plant tissues, consume high electrical energy and as fluorescent lamps do not emit the exact required light wavelengths for optimum plant growth
Scientific experiments:
A large number of plant species have been assessed in greenhouse trials to make sure plants have higher quality in biomass and biochemical ingredients even higher or comparable with field conditions. Plant performance of mint, basil, lentil, lettuce, cabbage, parsley, carrot and… were measured by assessing health and vigor of plants and success in promoting growth. Promoting in profuse flowering of select ornamentals including primula, marigold, stock and…… were also noticed. Experiments unraveled surprising performance and production of vegetables and ornamental plants under LED light sources. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Durham tube**
Durham tube:
Durham tubes are used in microbiology to detect production of gas by microorganisms. They are simply smaller test tubes inserted upside down in another test tube so they are freely movable. The culture media to be tested is then added to the larger tube and sterilized, which also eliminates the initial air gap produced when the tube is inserted upside down. The culture media typically contains a single substance to be tested with the organism, such as to determine whether an organism can ferment a particular carbohydrate. After inoculation and incubation, any gas that is produced will form a visible gas bubble inside the small tube. Litmus solution can also be added to the culture media to give a visual representation of pH changes that occur during the production of gas. The method was first reported in 1898 by British microbiologist Herbert Durham.One limitation of the Durham tube is that it does not allow for precise determination of the type of gas that is produced within the inner tube, or measurements of the quantity of gas produced. However, Durham argued that quantitive measurements are of limited value because of the culture solution will absorb some of the gas in unknown, variable proportions. Additionally, using Durham tubes to provide evidence of fermentation may not be able to detect slow- or weakly-fermenting organisms when the resultant carbon dioxide diffuses back into the solution as quickly as it is formed, so a negative test using Durham tubes does not indicate decisive physiological significance. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**C-clamp**
C-clamp:
A C-clamp or G-clamp or G-cramp is a type of clamp device typically used to hold a wood or metal workpiece, and often used in, but are not limited to, carpentry and welding. Often believed that these clamps are called "C" clamps because of their C-shaped frame, or also often called C-clamps or G-clamps because including the screw part, they are shaped like an uppercase letter G. However, in fact, they were originally called a carriage maker's clamp, or Carriage Clamp.
Description:
C-clamps or G-clamps are typically made of steel or cast iron, though smaller clamps may be made of pot metal. At the top of the "C," is usually a small flat edge. At the bottom is a threaded hole through which a large threaded screw protrudes. One end of this screw contains a flat edge of similar size to the one at the top of the frame, and the other end usually a small metal bar, perpendicular to the screw itself, which is used to gain leverage when tightening the clamp. When the clamp is completely closed, the flat end of the screw is in contact with the flat end on the frame. When the clamp is actually used, it is very rare that this occurs. Generally, some other object or objects will be contained between the top and bottom flat edges.
Usage:
A G-clamp is used by means of turning the screw through the bottom of the frame until the desired state of pressure or release is reached. In the case that the clamp is being tightened, this is when the objects being secured are satisfactorily secured between the flat end of the screw and the flat end of the frame. If the clamp is being loosened, this is when a sufficient amount of force is released to allow the secured objects to be moved.
Usage:
Woodworking While a G-clamp is a useful tool for woodworking, special care should be taken when working with any woods. The flat gripping edges of the frame, generally no larger than half an inch or a centimeter (depending on the size of the clamp) can cause indentations and marring of the surfaces being clamped. This can be avoided by buffering between the clamp and the timber using two pieces of scrap wood. As each piece of scrap wood is directly in contact with the flat edges of the frame and with the items being clamped, this allows the scrap wood to receive the damage from the clamping, while dispersing the clamping force across the piece of scrap wood into the clamped objects. Deep-throated clamps are also available and provide greater reach for smaller jobs.
Usage:
Stage Lighting see C-Clamp (stagecraft) C-clamps are frequently used to hang stage lighting instruments. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Music technology (electronic and digital)**
Music technology (electronic and digital):
Digital music technology encompasses digital instruments, computers, electronic effects units, software, or digital audio equipment by a performer, composer, sound engineer, DJ, or record producer to produce, perform or record music. The term refers to electronic devices, instruments, computer hardware, and software used in performance, playback, recording, composition, mixing, analysis, and editing of music.
Education:
Professional training Courses in music technology are offered at many different Universities as part of degree programs focusing on performance, composition, music research at the undergraduate and graduate level. The study of music technology is usually concerned with the creative use of technology for creating new sounds, performing, recording, programming sequencers or other music-related electronic devices, and manipulating, mixing and reproducing music. Music technology programs train students for careers in "...sound engineering, computer music, audio-visual production and post-production, mastering, scoring for film and multimedia, audio for games, software development, and multimedia production." Those wishing to develop new music technologies often train to become an audio engineer working in R&D. Due to the increasing role of interdisciplinary work in music technology, individuals developing new music technologies may also have backgrounds or training in computer programming, computer hardware design, acoustics, record producing or other fields.
Education:
Use of music technology in education Digital music technologies are widely used to assist in music education for training students in the home, elementary school, middle school, high school, college and university music programs. Electronic keyboard labs are used for cost-effective beginner group piano instruction in high schools, colleges, and universities. Courses in music notation software and basic manipulation of audio and MIDI can be part of a student's core requirements for a music degree. Mobile and desktop applications are available to aid the study of music theory and ear training. Digital pianos, such as those offered by Roland, provide interactive lessons and games using the built-in features of the instrument to teach music fundamentals.
History:
Development of digital musical technologies can be traced back to the analog music technologies of the early 20th century, such as the electromechanical Hammond organ, which was invented in 1929. In the 2010s, the ontological range of music technology has greatly increased, and it may now be electronic, digital, software-based or indeed even purely conceptual.Early pioneers included Luigi Russolo, Halim El-Dabh, Pierre Schaeffer, Pierre Henry, Edgard Varèse, Karlheinz Stockhausen, Ikutaro Kakehashi, King Tubby., and others who manipulated sounds using tape machines—splicing tape and changing its playback speed to alter pre-recorded samples. Pierre Schaefer was credited for inventing this method of composition, known as musique concrète, in 1948 in Paris, France. In this style of composition, existing material is manipulated to create new timbres. Musique concrète contrasts a later style that emerged in the mid-1950s in Cologne, Germany, known as elektronische Musik. This style, invented by Karlheinz Stockhausen, involves creating new sounds without the use of pre-existing material. Unlike musique concrète, which primarily focuses on timbre, elektronische Musik focuses on structure. Influences of these two styles still prevail today in today's modern music and music technology. The concept of the software digital audio workstation is the emulation of a traditional recording studio. Colored strips, known as regions, can be spliced, stretched, and re-ordered, analogous to tape. Similarly, software representations of classic synthesizers emulate their analog counterparts.
History:
Analog synthesizer history Classic analog synthesizers include the Moog Minimoog, ARP Odyssey, Yamaha CS-80, Korg MS-20, Sequential Circuits Prophet-5, Roland TB-303, Roland Alpha Juno. The most iconic bass synthesizer is the Roland TB-303, widely used in acid house music.
Digital synthesizer history Classic digital synthesizers include the Fairlight CMI, PPG Wave, Nord Modular and Korg M1.
History:
Digital synthesizer in Japan Through the 1970s and 1980s, Japanese synthesizer manufacturers produced more affordable synthesizers than those produced in America, with synthesizers made by Yamaha Corporation, Roland Corporation, Korg, Kawai and other companies. Yamaha's DX7 was one of the first mass-market, relatively inexpensive synthesizer keyboards. The DX7 is an FM synthesis based digital synthesizer manufactured from 1983 to 1989. It was the first commercially successful digital synthesizer. Its distinctive sound can be heard on many recordings, especially pop music from the 1980s. The monotimbral, 16-note polyphonic DX7 was the moderately priced model of the DX series keyboard synthesizers. Over 200,000 of the original DX7 were made, and it remains one of the best-selling synthesizers of all time.
History:
Computer music history Max Mathews Computer and synthesizer technology joining together changed the way music is made, and is one of the fastest changing aspects of music technology today. Max Mathews, a telecommunications engineer at Bell Telephone Laboratories' Acoustic and Behavioural Research Department, is responsible for some of the first digital music technology in the 50s. Max Mathews also pioneered a cornerstone of music technology; analog to digital conversion.At Bell Laboratories, Matthews conducted research to improve the telecommunications quality for long-distance phone calls. Owing to long-distance and low-bandwidth, audio quality over phone calls across the United States was poor. Thus, Matthews devised a method in which sound was synthesized via computer on the distant end rather than transmitted. Matthews was an amateur violinist, and during a conversation with his superior, John Pierce at Bell Labs, Pierce posed the idea of synthesizing music through a computer since Matthews had already synthesized speech. Matthews agreed, and beginning in the 1950s wrote a series of programs known as MUSIC. MUSIC consisted of two files—and orchestra file containing data telling the computer how to synthesize sound—and a score file instructing the program what notes to play using the instruments defined in the orchestra file. Matthews wrote five iterations of MUSIC, calling them MUSIC I-V respectively. Subsequently, as the program was adapted and expanded as it was written to run on various platforms, its name changed to reflect its new changes. This series of programs became known as the MUSICn paradigm. The concept of the MUSIC now exists in the form of Csound.Later Max Matthews worked as an advisor to IRCAM (Institut de recherche et coordination acoustique/musique; English: Institute for Research and Coordination in Acoustics/Music) in the late 1980s. There, he taught Miller Puckette, a researcher. Puckette developed a program in which music could be programmed graphically. The program could transmit and receive MIDI messages to generate interactive music in real-time. Inspired by Matthews, Puckette named the program Max. Later, a researcher named David Zicarelli visited IRCAM, saw the capabilities of Max and felt it could be developed further. He took a copy of Max with him when he left and eventually added capabilities to process audio signals. Zicarelli named this new part of the program MSP after Miller Puckette. Zicarelli developed the commercial version of MaxMSP and sold it at his company, Cycling '74, beginning in 1997. The company has since been acquired by Ableton.
History:
Later history The first generation of professional commercially available computer music instruments, or workstations as some companies later called them, were very sophisticated elaborate systems that cost a great deal of money when they first appeared. They ranged from $25,000 to $200,000. The two most popular were the Fairlight, and the Synclavier.
History:
It was not until the advent of MIDI that general-purpose computers started to play a role in music production. Following the widespread adoption of MIDI, computer-based MIDI editors and sequencers were developed. MIDI-to-CV/Gate converters were then used to enable analogue synthesizers to be controlled by a MIDI sequencer.Reduced prices in personal computers caused the masses to turn away from the more expensive workstations. Advancements in technology have increased the speed of hardware processing and the capacity of memory units. Powerful programs for sequencing, recording, notating, and mastering music.
History:
MIDI history At the NAMM Show of 1983 in Los Angeles, MIDI was released. A demonstration at the convention showed two previously incompatible analog synthesizers, the Prophet 600 and Roland Jupiter-6, communicating with each other, enabling a player to play one keyboard while getting the output from both of them. This was a massive breakthrough in the 1980s, as it allowed synths to be accurately layered in live shows and studio recordings. MIDI enables different electronic instruments and electronic music devices to communicate with each other and with computers. The advent of MIDI spurred a rapid expansion of the sales and production of electronic instruments and music software.
History:
In 1985, several of the top keyboard manufacturers created the MIDI Manufacturers Association (MMA). This newly founded association standardized the MIDI protocol by generating and disseminating all the documents about it. With the development of the MIDI File Format Specification by Opcode, every music software company's MIDI sequencer software could read and write each other's files.
Since the 1980s, personal computers developed and became the ideal system for utilizing the vast potential of MIDI. This has created a large consumer market for software such as MIDI-equipped electronic keyboards, MIDI sequencers and digital audio workstations. With universal MIDI protocols, electronic keyboards, sequencers, and drum machines can all be connected together.
Vocal synthesis history until 1980s VODER on Bell Lab.
History:
Coinciding with the history of computer music is the history of vocal synthesis. Prior to Max Matthews synthesizing speech with a computer, analog devices were used to recreate speech. In the 1930s, an engineer named Holmer Dudley invented the VODER (Voice Operated Demonstrator), an electro-mechanical device which generated a sawtooth wave and white-noise. Various parts of the frequency spectrum of the waveforms could be filtered to generate the sounds of speech. Pitch was modulated via a bar on a wrist strap worn by the operator. In the 1940s Dudley, invented the VOCODER (Voice Operated Coder). Rather than synthesizing speech from scratch, this machine operated by accepting incoming speech and breaking it into its spectral components. In the late 1960s and early 1970s, bands and solo artists began using the VOCODER to blend speech with notes played on a synthesizer.
History:
Singing Kelly-Lochbaum Vocal Tract on Bell Lab.
History:
Meanwhile, at Bell Laboratories, Max Matthews worked with researchers Kelly and Lachbaum to develop a model of the vocal tract to study how its prosperities contributed to speech generation. Using the model of the vocal tract, Matthews used linear predictive coding (LPC)—a method in which a computer estimates the formants and spectral content of each word based on information about the vocal model, including various applied filters representing the vocal tract—to make a computer (an IBM 704) sing for the first time in 1962. The computer performed a rendition of "Bicycle Built for Two." CHANT on IRCAM In the 1970s at IRCAM in France, researchers developed a piece of software called CHANT (French for "sing"). CHANT was based FOF (Fomant ond Formatique) synthesis, in which the peak frequencies of a sound are created and shaped using granular synthesis—as opposed to filtering frequencies to create speech.
History:
Concatenation synthesis using MIDI Through the 1980s and 1990s as MIDI devices became commercially available, speech was generated by mapping MIDI data to samples of the components of speech stored in sample libraries.
Synthesizers and drum machines:
Synthesizers A synthesizer is an electronic musical instrument that generates electric signals that are converted to sound through instrument amplifiers and loudspeakers or headphones. Synthesizers may either imitate existing sounds (instruments, vocal, natural sounds, etc.), or generate new electronic timbres or sounds that did not exist before. They are often played with an electronic musical keyboard, but they can be controlled via a variety of other input devices, including music sequencers, instrument controllers, fingerboards, guitar synthesizers, wind controllers, and electronic drums. Synthesizers without built-in controllers are often called sound modules, and are controlled using a controller device.
Synthesizers and drum machines:
Synthesizers use various methods to generate a signal. Among the most popular waveform synthesis techniques are subtractive synthesis, additive synthesis, wavetable synthesis, frequency modulation synthesis, phase distortion synthesis, physical modeling synthesis and sample-based synthesis. Other less common synthesis types include subharmonic synthesis, a form of additive synthesis via subharmonics (used by mixture trautonium), and granular synthesis, sample-based synthesis based on grains of sound, generally resulting in soundscapes or clouds. In the 2010s, synthesizers are used in many genres of pop, rock and dance music. Contemporary classical music composers from the 20th and 21st century write compositions for synthesizer.
Synthesizers and drum machines:
Drum machines A drum machine is an electronic musical instrument designed to imitate the sound of drums, cymbals, other percussion instruments, and often basslines. Drum machines either play back prerecorded samples of drums and cymbals or synthesized re-creations of drum/cymbal sounds in a rhythm and tempo that is programmed by a musician. Drum machines are most commonly associated with electronic dance music genres such as house music, but are also used in many other genres. They are also used when session drummers are not available or if the production cannot afford the cost of a professional drummer. In the 2010s, most modern drum machines are sequencers with a sample playback (rompler) or synthesizer component that specializes in the reproduction of drum timbres. Though features vary from model to model, many modern drum machines can also produce unique sounds, and allow the user to compose unique drum beats and patterns.
Synthesizers and drum machines:
Electro-mechanical drum machines were first developed in 1949, with the invention of the Chamberlin Rhythmate. Transistorized electronic drum machines Seeburg Select-A-Rhythm appeared in 1964.Classic drum machines include the Korg Mini Pops 120, PAiA Programmable Drum Set, Roland CR-78, LinnDrum, Roland TR-909, Oberheim DMX, E-MU SP-12, Alesis HR-16, and Elektron SPS1 Machinedrum (in chronological order).
Synthesizers and drum machines:
Drum machines in Japan The Ace Tone Rhythm Ace, created by Ikutaro Kakehashi, began appearing in popular music from the late 1960s, followed by drum machines from Korg and Ikutaro's later Roland Corporation also appearing in popular music from the early 1970s. Sly and the Family Stone's 1971 album There's a Riot Goin' On helped to popularize the sound of early drum machines, along with Timmy Thomas' 1972 R&B hit "Why Can't We Live Together" and George McCrae's 1974 disco hit "Rock Your Baby" which used early Roland rhythm machines.Early drum machines sounded drastically different than the drum machines that gained their peak popularity in the 1980s and defined an entire decade of pop music. The most iconic drum machine was the Roland TR-808, widely used in hip hop and dance music.
Sampling technology after 1980s:
Digital sampling technology, introduced in the 1970s, has become a staple of music production in the 2000s. Devices that use sampling, record a sound digitally (often a musical instrument, such as a piano or flute being played), and replay it when a key or pad on a controller device (e.g., an electronic keyboard, electronic drum pad, etc.) is pressed or triggered. Samplers can alter the sound using various audio effects and audio processing. Sampling has its roots in France with the sound experiments carried out by musique concrète practitioners.
Sampling technology after 1980s:
In the 1980s, when the technology was still in its infancy, digital samplers cost tens of thousands of dollars and they were only used by the top recording studios and musicians. These were out of the price range of most musicians. Early samplers include the 8-bit Electronic Music Studios MUSYS-3 circa 1970, Computer Music Melodian in 1976, Fairlight CMI in 1979, Emulator I in 1981, Synclavier II Sample-to-Memory (STM) option circa 1980, Ensoniq Mirage in 1984, and Akai S612 in 1985. The latter's successor, the Emulator II (released in 1984), listed for $8,000. Samplers were released during this period with high price tags, such as the K2000 and K2500.
Sampling technology after 1980s:
Some important hardware samplers include the Kurzweil K250, Akai MPC60, Ensoniq Mirage, Ensoniq ASR-10, Akai S1000, E-mu Emulator, and Fairlight CMI.One of the biggest uses of sampling technology was by hip-hop music DJs and performers in the 1980s. Before affordable sampling technology was readily available, DJs would use a technique pioneered by Grandmaster Flash to manually repeat certain parts in a song by juggling between two separate turntables. This can be considered as an early precursor of sampling. In turn, this turntablism technique originates from Jamaican dub music in the 1960s, and was introduced to American hip hop in the 1970s.
Sampling technology after 1980s:
In the 2000s, most professional recording studios use digital technologies. In recent years, many samplers have only included digital technology. This new generation of digital samplers are capable of reproducing and manipulating sounds. Digital sampling plays an integral part in some genres of music, such as hip-hop and trap. Advanced sample libraries have made complete performances of orchestral compositions possible that sound similar to a live performance. Modern sound libraries allow musicians to have the ability to use the sounds of almost any instrument in their productions.
Sampling technology after 1980s:
Sampling technology in Japan Early samplers include the 12-bit Toshiba LMD-649 in 1981.The first affordable sampler in Japan was the Ensoniq Mirage in 1984. Also the AKAI S612 became available in 1985, retailed for US$895. Other companies soon released affordable samplers, including Oberheim DPX-1 in 1987, and more by Korg, Casio, Yamaha, and Roland. Some important hardware samplers in Japan include the Akai Z4/Z8, Roland V-Synth, Casio FZ-1.
MIDI:
MIDI has been the musical instrument industry standard interface since the 1980s through to the present day. It dates back to June 1981, when Roland Corporation founder Ikutaro Kakehashi proposed the concept of standardization between different manufacturers' instruments as well as computers, to Oberheim Electronics founder Tom Oberheim and Sequential Circuits president Dave Smith. In October 1981, Kakehashi, Oberheim and Smith discussed the concept with representatives from Yamaha, Korg and Kawai. In 1983, the MIDI standard was unveiled by Kakehashi and Smith.Some universally accepted varieties of MIDI software applications include music instruction software, MIDI sequencing software, music notation software, hard disk recording/editing software, patch editor/sound library software, computer-assisted composition software, and virtual instruments. Current developments in computer hardware and specialized software continue to expand MIDI applications.
Computers in music technology after 1980s:
Following the widespread adoption of MIDI, computer-based MIDI editors and sequencers were developed. MIDI-to-CV/Gate converters were then used to enable analogue synthesizers to be controlled by a MIDI sequencer.Reduced prices in personal computers caused the masses to turn away from the more expensive workstations. Advancements in technology have increased the speed of hardware processing and the capacity of memory units. Software developers write new, more powerful programs for sequencing, recording, notating, and mastering music.
Computers in music technology after 1980s:
Digital audio workstation software, such as Pro Tools, Logic, and many others, have gained popularity among the vast array of contemporary music technology in recent years. Such programs allow the user to record acoustic sounds with a microphone or software instrument, which may then be layered and organized along a timeline and edited on a flat-panel display of a computer. Recorded segments can be copied and duplicated ad infinitum, without any loss of fidelity or added noise (a major contrast from analog recording, in which every copy leads to a loss of fidelity and added noise). Digital music can be edited and processed using a multitude of audio effects. Contemporary classical music sometimes uses computer-generated sounds—either pre-recorded or generated and manipulated live—in conjunction or juxtaposed on classical acoustic instruments like the cello or violin. Music is scored with commercially available notation software.In addition to the digital audio workstations and music notation software, which facilitate the creation of fixed media (material that does not change each time it is performed), software facilitating interactive or generative music continues to emerge. Composition based on conditions or rules (algorithmic composition) has given rise to software which can automatically generate music based on input conditions or rules. Thus, the resulting music evolves each time conditions change. Examples of this technology include software designed for writing music for video games—where music evolves as a player advances through a level or when certain characters appear—or music generated from artificial intelligence trained to convert biometrics like EEG or ECG readings into music. Because this music is based on user interaction, it will be different each time it is heard. Other examples of generative music technology include the use of sensors connected to computer and artificial intelligence to generate music based on captured data, such as environmental factors, the movements of dancers, or physical inputs from a digital device such as a mouse or game controller. Software applications offering capabilities for generative and interactive music include SuperCollider, MaxMSP/Jitter, and Processing. Interactive music is made possible through physical computing, where the data from the physical world affects a computer's output and vice versa.
Singing synthesis after 2010s:
In the 2010s, Singing synthesis technology has taken advantage of the recent advances in artificial intelligence—deep listening and machine learning to better represent the nuances of the human voice. New high fidelity sample libraries combined with digital audio workstations facilitate editing in fine detail, such as shifting of formats, adjustment of vibrato, and adjustments to vowels and consonants. Sample libraries for various languages and various accents are available. With today's advancements in Singing synthesis, artists sometimes use sample libraries in lieu of backing singers.
Timeline:
1917 : Leon Theremin invented the prototype of the Theremin 1944 : Halim El-Dabh produces earliest electroacoustic tape music 1952 : Harry F. Olson and Herbert Belar invent the RCA Synthesizer 1952 : Osmand Kendal develops the Composer-Tron for the Marconi Wireless Company 1956 : Raymond Scott develops the Clavivox 1958 : Yevgeny Murzin along with several colleagues create the ANS synthesizer 1959 : Wurlitzer manufactures The Sideman, the first commercial electro-mechanical drum machine 1963 : The Mellotron starts to be manufactured in London 1964 : The Moog synthesizer is released 1968 : King Tubby pioneers dub music, an early form of popular electronic music 1970 : ARP 2600 is manufactured 1982 : Sony and Philips introduce compact disc 1983 : Introduction of MIDI 1986 : The first digital consoles appear 1987 : Digidesign markets Sound Tools Timeline in Japan 1963 : Keio Electronics (later Korg) produces the DA-20 1964 : Ikutaro Kakehashi debuts Ace Tone R-1 Rhythm Ace, the first electronic drum 1965 : Nippon Columbia patents an early electronic drum machine 1966 : Korg releases Donca-Matic DE-20, an early electronic drum machine 1967 : Ace Tone releases FR-1 Rhythm Ace, the first drum machine to enter popular music 1967 : First PCM recorder developed by NHK 1969 : Matsushita engineer Shuichi Obata invents first direct-drive turntable, Technics SP-10 1973 : Yamaha release Yamaha GX-1, the first polyphonic synthesizer 1974 : Yamaha build first digital synthesizer 1977 : Roland release MC-8, an early microprocessor-driven CV/Gate digital sequencer 1978 : Roland releases CR-78, the first microprocessor-driven drum machine 1979 : Casio releases VL-1, the first commercial digital synthesizer 1980 : Roland releases TR-808, the most widely used drum machine in popular music 1980 : Roland introduces DCB protocol and DIN interface with TR-808 1980 : Yamaha releases GS-1, the first FM digital synthesizer 1980 : Kazuo Morioka creates Firstman SQ-01, the first bass synth with a sequencer 1981 : Roland releases TB-303, a bass synthesizer that lays foundations for acid house music 1981 : Toshiba's LMD-649, the first PCM digital sampler in Japan, introduced with Yellow Magic Orchestra's Technodelic 1982 : First MIDI synthesizers released, Roland Jupiter-6 and Prophet 600 1983 : Roland releases MSQ-700, the first MIDI sequencer 1983 : Roland releases TR-909, the first MIDI drum machine 1983 : Yamaha releases DX7, the first commercially successful digital synthesizer 1985 : Akai releases the Akai S612, a digital sampler 1988 : Akai introduces the Music Production Controller (MPC) series of digital samplers 1994 : Yamaha unveils the ProMix 01 | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Linear B Syllabary**
Linear B Syllabary:
Linear B Syllabary is a Unicode block containing characters for the syllabic writing of Mycenaean Greek.
History:
The following Unicode-related documents record the purpose and process of defining specific characters in the Linear B Syllabary block: | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Verbosus**
Verbosus:
Verbosus is a browser-based LaTeX editor which allows a user to create and handle LaTeX projects in a browser. The graphical user interface (GUI) does deliberately not resemble non-browser-based Editors such as TeXworks. It was designed to function and being used in a browser.
Verbosus requires no installation of any software packages like MiKTeX, TeX Live, etc.
As it is the case in other non-browser-based Latex tools a PDF-Viewer is integrated which allows to generate a .pdf out of the Latex-Code. A preview is displayed on the side of the Latex Code. Additionally, the editor supports syntax highlighting which increases the readability as well as code completion.
The connection between the web browser and server is secured by using the HTTPS protocol which provides encryption and secure identification of the server.
Technology:
The client side of Verbosus was developed entirely in JavaScript. In addition, it uses the Dojo Toolkit for server communication. The integrated approach of using JavaScript-On-Demand allows the (for the user transparent) reloading of contents without the need to reload the whole page. Since JavaScript is used as the core-technology no additional plug-in (like Adobe Flash Player, etc.) is required.
Mobile platforms:
VerbTeX is an application for Android, iOS and Windows that uses the interface of Verbosus to synchronize data and generate the PDF from LaTeX code. Since version 3.4 there exists a native Android App called Anoc which is an Octave editor that also uses the interface of Verbosus to generate the result and create plots.
Collaboration:
Verbosus supports simultaneous, collaborative editing of the same project by different users. A built-in tool that implements the diff-algorithm displays conflicting changes and allows a user to resolve conflicts after deciding which version is the correct one. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Khmaladze transformation**
Khmaladze transformation:
In statistics, the Khmaladze transformation is a mathematical tool used in constructing convenient goodness of fit tests for hypothetical distribution functions. More precisely, suppose X1,…,Xn are i.i.d., possibly multi-dimensional, random observations generated from an unknown probability distribution. A classical problem in statistics is to decide how well a given hypothetical distribution function F , or a given hypothetical parametric family of distribution functions {Fθ:θ∈Θ} , fits the set of observations. The Khmaladze transformation allows us to construct goodness of fit tests with desirable properties. It is named after Estate V. Khmaladze.
Khmaladze transformation:
Consider the sequence of empirical distribution functions Fn based on a sequence of i.i.d random variables, X1,…,Xn , as n increases. Suppose F is the hypothetical distribution function of each Xi . To test whether the choice of F is correct or not, statisticians use the normalized difference, vn(x)=n[Fn(x)−F(x)].
Khmaladze transformation:
This vn , as a random process in x , is called the empirical process. Various functionals of vn are used as test statistics. The change of the variable vn(x)=un(t) , t=F(x) transforms to the so-called uniform empirical process un . The latter is an empirical processes based on independent random variables Ui=F(Xi) , which are uniformly distributed on [0,1] if the Xi s do indeed have distribution function F This fact was discovered and first utilized by Kolmogorov (1933), Wald and Wolfowitz (1936) and Smirnov (1937) and, especially after Doob (1949) and Anderson and Darling (1952), it led to the standard rule to choose test statistics based on vn . That is, test statistics ψ(vn,F) are defined (which possibly depend on the F being tested) in such a way that there exists another statistic φ(un) derived from the uniform empirical process, such that ψ(vn,F)=φ(un) . Examples are sup sup sup sup t|un(t)|a(t) and ∫−∞∞vn2(x)dF(x)=∫01un2(t)dt.
Khmaladze transformation:
For all such functionals, their null distribution (under the hypothetical F ) does not depend on F , and can be calculated once and then used to test any F However, it is only rarely that one needs to test a simple hypothesis, when a fixed F as a hypothesis is given. Much more often, one needs to verify parametric hypotheses where the hypothetical F=Fθn , depends on some parameters θn , which the hypothesis does not specify and which have to be estimated from the sample X1,…,Xn itself.
Khmaladze transformation:
Although the estimators θ^n , most commonly converge to true value of θ , it was discovered that the parametric, or estimated, empirical process v^n(x)=n[Fn(x)−Fθ^n(x)] differs significantly from vn and that the transformed process u^n(t)=v^n(x) , t=Fθ^n(x) has a distribution for which the limit distribution, as n→∞ , is dependent on the parametric form of Fθ and on the particular estimator θ^n and, in general, within one parametric family, on the value of θ From mid-1950s to the late-1980s, much work was done to clarify the situation and understand the nature of the process v^n In 1981, and then 1987 and 1993, Khmaladze suggested to replace the parametric empirical process v^n by its martingale part wn only.
Khmaladze transformation:
v^n(x)−Kn(x)=wn(x) where Kn(x) is the compensator of v^n(x) . Then the following properties of wn were established: Although the form of Kn , and therefore, of wn , depends on Fθ^n(x) , as a function of both x and θn , the limit distribution of the time transformed process ωn(t)=wn(x),t=Fθ^n(x) is that of standard Brownian motion on [0,1] , i.e., is again standard and independent of the choice of Fθ^n .The relationship between v^n and wn and between their limits, is one to one, so that the statistical inference based on v^n or on wn are equivalent, and in wn , nothing is lost compared to v^n The construction of innovation martingale wn could be carried over to the case of vector-valued X1,…,Xn , giving rise to the definition of the so-called scanning martingales in Rd .For a long time the transformation was, although known, still not used. Later, the work of researchers like Koenker, Stute, Bai, Koul, Koening, and others made it popular in econometrics and other fields of statistics. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Thermoreceptor**
Thermoreceptor:
A thermoreceptor is a non-specialised sense receptor, or more accurately the receptive portion of a sensory neuron, that codes absolute and relative changes in temperature, primarily within the innocuous range. In the mammalian peripheral nervous system, warmth receptors are thought to be unmyelinated C-fibres (low conduction velocity), while those responding to cold have both C-fibers and thinly myelinated A delta fibers (faster conduction velocity). The adequate stimulus for a warm receptor is warming, which results in an increase in their action potential discharge rate. Cooling results in a decrease in warm receptor discharge rate. For cold receptors their firing rate increases during cooling and decreases during warming. Some cold receptors also respond with a brief action potential discharge to high temperatures, i.e. typically above 45 °C, and this is known as a paradoxical response to heat. The mechanism responsible for this behavior has not been determined.
Location:
In humans, along the axons of Lissauer's tract temperature or pressure sensations enter the spinal cord. The Lissauer's tract will synapse on first-order neurons in grey matter of the dorsal horn, one or two vertebral levels up. The axons of these second-order neurons then decussate, joining the spinothalamic tract as they ascend to neurons in the ventral posterolateral nucleus of the thalamus.
Location:
In mammals, temperature receptors innervate various tissues including the skin (as cutaneous receptors), cornea and urinary bladder. Neurons from the pre-optic and hypothalamic regions of the brain that respond to small changes in temperature have also been described, providing information on core temperature. The hypothalamus is involved in thermoregulation, the thermoreceptors allowing feed-forward responses to a predicted change in core body temperature in response to changing environmental conditions.
Structure:
Thermoreceptors have been classically described as having 'free' non-specialized endings; the mechanism of activation in response to temperature changes is not completely understood.
Function:
Cold-sensitive thermoreceptors give rise to the sensations of cooling, cold and freshness.
In the cornea cold receptors are thought to respond with an increase in firing rate to cooling produced by evaporation of lacrimal fluid 'tears' and thereby to elicit a blink reflex.
Function:
Other thermoreceptors will react to opposite triggers and give rise to heat and in some cases even burning sensations. This is often experienced when coming in contact with capsaicin, an active chemical commonly found in chili peppers. When coming in contact with your tongue (or any internal surface), the capsaicin de-polarizes the nerve fibers, allowing sodium and calcium into the fibers. In order for fibers to do so, they must have a specific thermoreceptor. The thermoreceptor reacting to capsaicin and other heat producing chemicals is known as TRPV1. In response to heat, the TRPV1 receptor opens up passages that allow ions to pass through, causing the sensation of heat or burning.
Function:
TRPV1 also has a molecular cousin, TRPM8. Unlike TRPV1, TRPM8 produces cooling sensations as mentioned previously. Similar to TRPV1, TRPM8 responds to a certain chemical trigger by opening its ion pathways. In this case, the chemical trigger is often menthol or other cooling agents. Studies performed on mice determined that the presence of both these receptors allows for a gradient of temperature sensing. Mice lacking the TRPV1 receptor were still capable of determining areas significantly colder than on a heated platform. Mice lacking the TRPM8 receptor however, were not able to determine the difference between a warm platform and a cold platform, suggesting we rely on TRPM8 to determine cold feelings and sensations.
Distribution:
Warm and cold receptors play a part in sensing innocuous environmental temperature. Temperatures likely to damage an organism are sensed by sub-categories of nociceptors that may respond to noxious cold, noxious heat or more than one noxious stimulus modality (i.e., they are polymodal). The nerve endings of sensory neurons that respond preferentially to cooling are found in moderate density in the skin but also occur in relatively high spatial density in the cornea, tongue, bladder, and facial skin. The speculation is that lingual cold receptors deliver information that modulates the sense of taste; i.e. some foods taste good when cold, while others do not.
Mechanism of transduction:
This area of research has recently received considerable attention with the identification and cloning of the Transient Receptor Potential (TRP) family of proteins. The transduction of temperature in cold receptors is mediated in part by the TRPM8 channel. This channel passes a mixed inward cationic (predominantly carried by Na+ ions although the channel is also permeable to Ca2+) current of a magnitude that is inversely proportional to temperature. The channel is sensitive over a temperature range spanning about 10-35 °C. TRPM8 can also be activated by the binding of an extracellular ligand. Menthol can activate the TRPM8 channel in this way. Since the TRPM8 is expressed in neurons whose physiological role is to signal cooling, menthol applied to various bodily surfaces evokes a sensation of cooling. The feeling of freshness associated with the activation of cold receptors by menthol, particularly those in facial areas with axons in the trigeminal (V) nerve, accounts for its use in numerous toiletries including toothpaste, shaving lotions, facial creams and the like.
Mechanism of transduction:
Another molecular component of cold transduction is the temperature dependence of so-called leak channels which pass an outward current carried by potassium ions. Some leak channels derive from the family of two-pore (2P) domain potassium channels. Amongst the various members of the 2P-domain channels, some close quite promptly at temperatures less than about 28 °C (e.g. KCNK4(TRAAK), TREK). Temperature also modulates the activity of the Na+/K+-ATPase. The Na+/K+-ATPase is a P-type pump that extrudes 3Na+ ions in exchange for 2K+ ions for each hydrolytic cleavage of ATP. This results in a net movement of positive charge out of the cell, i.e. a hyperpolarizing current. The magnitude of this current is proportional to the rate of pump activity.
Mechanism of transduction:
It has been suggested that it is the constellation of various thermally sensitive proteins together in a neuron that gives rise to a cold receptor. This emergent property of the neuron is thought to comprise, the expression of the aforementioned proteins as well as various voltage-sensitive channels including the hyperpolarization-activated, cyclic nucleotide-gated (HCN) channel and the rapidly activating and inactivating transient potassium channel (IKA). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Software as a service**
Software as a service:
Software as a service (SaaS ) is a software licensing and delivery model in which software is licensed on a subscription basis and is centrally hosted. SaaS is also known as on-demand software, web-based software, or web-hosted software.SaaS is considered to be part of cloud computing, along with several other as a service business models. SaaS apps are typically accessed by users of a web browser (a thin client). SaaS became a common delivery model for many business applications, including office software, messaging software, payroll processing software, DBMS software, management software, CAD software, development software, gamification, virtualization, accounting, collaboration, customer relationship management (CRM), management information systems (MIS), enterprise resource planning (ERP), invoicing, field service management, human resource management (HRM), talent acquisition, learning management systems, content management (CM), geographic information systems (GIS), and service desk management.
Software as a service:
SaaS has been incorporated into the strategies of nearly all enterprise software companies.
History:
Centralized hosting of business applications dates back to the 1960s. Starting in that decade, IBM and other mainframe computer providers conducted a service bureau business, often referred to as time-sharing or utility computing. Such services included offering computing power and database storage to banks and other large organizations from their worldwide data centers.The expansion of the Internet during the 1990s brought about a new class of centralized computing, called application service providers (ASP). ASPs provided businesses with the service of hosting and managing specialized business applications to reduce costs through central administration and the provider's specialization in a particular business application. Two of the largest ASPs were USI, which was headquartered in the Washington, D.C., area, and Futurelink Corporation, headquartered in Irvine, California.Software as a service essentially extends the idea of the ASP model. The term software as a service (SaaS), however, is commonly used in more specific settings: While most initial ASPs focused on managing and hosting third-party independent software vendors' software, contemporary SaaS offerings are typically provided by the software developer.
History:
Whereas many initial ASPs offered more traditional client-server applications, which require the installation of software on users' personal computers, later implementations can be Web applications that only require a web browser to use.
History:
Whereas the software architecture used by most initial ASPs mandated maintaining a separate instance of the application for each business, SaaS services can utilize a multi-tenant architecture, in which the application serves multiple businesses and users, and partitions its data accordingly. The acronym first appeared in the goods and services description of a USPTO trademark, filed on September 23, 1985. DbaaS (Database as a service) has emerged as a sub-variety of SaaS and is a type of cloud database.Microsoft referred to SaaS as "software plus services" for a few years.
Distribution and pricing:
The cloud (or SaaS) model has no physical need for indirect distribution because it is not distributed physically and is deployed almost instantaneously, thereby negating the need for traditional partners and middlemen. Unlike traditional software, which is conventionally sold as a perpetual license with an up-front cost (and an optional ongoing support fee), SaaS providers generally price applications using a subscription fee, most commonly a monthly fee or an annual fee. Consequently, the initial setup cost for SaaS is typically lower than the equivalent enterprise software. SaaS vendors typically price their applications based on some usage parameters, such as the number of users using the application. However, because in a SaaS environment customers' data reside with the SaaS vendor, opportunities also exist to charge per transaction, event, or other units of value, such as the number of processors required.The relatively low cost for user provisioning (i.e., setting up a new customer) in a multi-tenant environment enables some SaaS vendors to offer applications using the freemium model. In this model, a free service is made available with limited functionality or scope, and fees are charged for enhanced functionality or larger scope.A key driver of SaaS growth is SaaS vendors' ability to provide a price that is competitive with on-premises software. This is consistent with the traditional rationale for outsourcing IT systems, which involves applying economies of scale to application operation, i.e., an outside service provider may be able to offer better, cheaper, more reliable applications.
Architecture:
Most SaaS providers offer a multi-tenant architecture. With this model, a single version of the application, with a single configuration (hardware, network, operating system), is used for all customers ("tenants"). To support scalability, the application can be installed on multiple machines (called horizontal scaling). In some cases, a second version of the application is set up to offer a select group of customers access to pre-release versions of the applications (e.g., a beta version) for testing purposes. This is contrasted with traditional software, where multiple physical copies of the software — each potentially of a different version, with a potentially different configuration, and often customized — are installed across various customer sites.Although an exception rather than the norm, some SaaS providers use mechanisms such as virtualization to manage a large number of customers in place of multitenancy. Whether multitenancy is a necessary component of software as a service is debatable.
Vertical vs horizontal SaaS:
Horizontal SaaS and vertical SaaS are different models of cloud computing services. Horizontal SaaS targets a broad variety of customers, generally without regard to their industry. Some popular examples of horizontal SaaS vendors are Salesforce and HubSpot. Vertical SaaS, on the other hand, refers to a niche market targeting a narrower variety of customers to meet their specific requirements.
Characteristics:
Although not all software-as-a-service applications share all the following traits, the characteristics below are common among many of them: Accelerated feature delivery SaaS applications are often updated more frequently than traditional software, in many cases on a weekly or monthly basis. This is enabled by several factors: The application is hosted centrally, so an update is decided and executed by the provider, not by customers.
Characteristics:
The application only has a single configuration, making development testing faster.
The application vendor does not have to expend resources updating and maintaining backdated software versions because there is only one version.
The application vendor has access to all customer data, expediting design and regression testing.
The service provider has access to user behavior within the application (usually via web analytics), making it easier to identify areas worthy of improvement.Accelerated feature delivery is further enabled by agile software development methodologies. Such methodologies, which evolved in the mid-1990s, provide a set of software development tools and practices to support frequent software releases.
Characteristics:
Open integration protocols SaaS applications predominantly offer integration protocols and application programming interfaces (APIs) that operate over a wide area network.The ubiquity of SaaS applications and other Internet services and the standardization of their API technology has spawned the development of mashups, which are lightweight applications that combine data, presentation, and functionality from multiple services, creating a compound service. Mashups further differentiate SaaS applications from on-premises software as the latter cannot be easily integrated outside a company's firewall.
Characteristics:
Collaborative (and "social") functionality Inspired by the development of different internet networking services and the so-called web 2.0 functionality, many SaaS applications offer features that let their users collaborate and share information.For example, many project management applications delivered in the SaaS model offer—in addition to traditional project planning functionality—collaboration features letting users comment on tasks and plans and share documents within and outside an organization. Several other SaaS applications let users vote on and offer new feature ideas.Although some collaboration-related functionality is also integrated into on-premises software, (implicit or explicit) collaboration between users or different customers is only possible with centrally hosted software.
Characteristics:
OpenSaaS OpenSaaS refers to software as a service (SaaS) based on open-source code. Like SaaS applications, Open SaaS is a web-based application hosted, supported, and maintained by a service provider. While the roadmap for Open SaaS applications is defined by its community of users, upgrades and product enhancements are managed by a central provider. The term was coined in 2011 by Dries Buytaert, creator of the Drupal content management framework.Andrew Hoppin, a former Chief Information Officer for the New York State Senate, has been a vocal advocate of OpenSaaS for government, calling it "the future of government innovation." He points to WordPress and Why Unified as a successful example of an OpenSaaS software delivery model that gives customers "the best of both worlds, and more options. The fact that it is open source means that they can start building their websites by self-hosting WordPress and customizing their website to their heart's content. Concurrently, the fact that WordPress is SaaS means that they don't have to manage the website at all -- they can simply pay WordPress.com to host it."
Adoption drivers:
Several important changes to the software market and technology landscape have facilitated the acceptance and growth of SaaS: The growing use of web-based user interfaces by applications, along with the proliferation of associated practices (e.g., web design), continuously decreased the need for traditional client-server applications. Consequently, traditional software vendor's investment in software based on fat clients has become a disadvantage (mandating ongoing support), opening the door for new software vendors' offering a user experience perceived as more "modern".
Adoption drivers:
The standardization of web page technologies (HTML, JavaScript, CSS), the increasing popularity of web development as a practice, and the introduction and ubiquity of web application frameworks like Ruby on Rails or Laravel (PHP) gradually reduced the cost of developing new software services and enabled new providers to challenge traditional vendors.
The increasing penetration of broadband Internet access enabled remote centrally hosted applications to offer speed comparable to on-premises software.
The standardization of the HTTPS protocol as part of the web stack provided universally available lightweight security that is sufficient for most everyday applications.
The introduction and wide acceptance of lightweight integration protocols such as Representational State Transfer (REST) and SOAP enabled affordable integration between SaaS applications (residing in the cloud) with internal applications over wide area networks and with other SaaS applications.
Adoption challenges:
Some limitations slow down the acceptance of SaaS and prohibit it from being used in some cases: Because data is stored on the vendor's servers rather than on-premises, data security becomes an issue, particularly with intellectual property.
SaaS applications are hosted in the cloud, far away from the application users. This introduces latency into the environment; for example, the SaaS model is not suitable for applications that demand response times in milliseconds (OLTP).
Multi-tenant architectures, which drive cost efficiency for service providers, limit customization of applications for large clients, inhibiting such applications from being used in scenarios (applicable mostly to large enterprises) for which such customization is necessary.
Some business applications require access to or integration with customers' current data. When such data are large in volume or sensitive (e.g. end-user's personal information), integrating them with remotely hosted software can be costly or risky or can conflict with data governance regulations.
Adoption challenges:
Constitutional search/seizure warrant laws do not protect all forms of SaaS dynamically stored data. The result is that a link is added to the chain of security where access to the data, and, by extension, misuse of these data, are limited only by the assumed honesty of third parties or government agencies able to access the data on their recognizance.
Adoption challenges:
Organizations that adopt SaaS may find they are forced into adopting new versions, which might result in unforeseen training costs, an increase in the probability that a user might make an error or instability from bugs in the newer software.
Should the vendor of the software go out of business or suddenly EOL the software, the user may lose access to their software unexpectedly, which could destabilize their organization's current and future projects, as well as leave the user with older data they can no longer access or modify.
Relying on an Internet connection means that data is transferred to and from a SaaS firm at Internet speeds rather than the potentially higher speeds of a firm's internal network.
Adoption challenges:
The Ability of the SaaS hosting company to guarantee the uptime level agreed in the SLA (Service Level Agreement) The reliance on SaaS applications and services can lead to SaaS sprawl within enterprises. These disparate applications and services can become challenging to maintain technically and administratively, leading to the proliferation of shadow IT. The standard model also has limitations: Compatibility with hardware, other software, and operating systems.
Adoption challenges:
Licensing and compliance problems (unauthorized copies of the software program putting the organization at risk of fines or litigation).
Maintenance, support, and patch revision processes.
Healthcare applications According to a survey by the Healthcare Information and Management Systems Society, 83% of US IT healthcare organizations are now using cloud services, with 9.3% planning to, whereas 67% of IT healthcare organizations are currently running SaaS-based applications.
Data escrow:
Software as a service data escrow is the process of keeping a copy of critical software-as-a-service application data with an independent third party. Similar to source code escrow, where critical software source code is stored with an independent third party, SaaS data escrow applies the same logic to the data within a SaaS application. It allows companies to protect and ensure all the data that resides within SaaS applications, protecting against data loss.There are many and varied reasons for considering SaaS data escrow, including concerns about vendor bankruptcy, unplanned service outages, and potential data loss or corruption.
Data escrow:
Many businesses either ensure that they are complying with their data governance standards or try to enhance their reporting and business analytics against their SaaS data.
Criticism:
One notable criticism of SaaS comes from Richard Stallman of the Free Software Foundation, who refers to it as Service as a Software Substitute (SaaSS). He considers the use of SaaSS to be a violation of the principles of free software. According to Stallman: With SaaSS, the users do not have even the executable file that does their computing: it is on someone else's server, where the users can't see or touch it. Thus it is impossible for them to ascertain what it really does, and impossible to change it. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Turing degree**
Turing degree:
In computer science and mathematical logic the Turing degree (named after Alan Turing) or degree of unsolvability of a set of natural numbers measures the level of algorithmic unsolvability of the set.
Overview:
The concept of Turing degree is fundamental in computability theory, where sets of natural numbers are often regarded as decision problems. The Turing degree of a set is a measure of how difficult it is to solve the decision problem associated with the set, that is, to determine whether an arbitrary number is in the given set.
Overview:
Two sets are Turing equivalent if they have the same level of unsolvability; each Turing degree is a collection of Turing equivalent sets, so that two sets are in different Turing degrees exactly when they are not Turing equivalent. Furthermore, the Turing degrees are partially ordered, so that if the Turing degree of a set X is less than the Turing degree of a set Y, then any (possibly noncomputable) procedure that correctly decides whether numbers are in Y can be effectively converted to a procedure that correctly decides whether numbers are in X. It is in this sense that the Turing degree of a set corresponds to its level of algorithmic unsolvability.
Overview:
The Turing degrees were introduced by Emil Leon Post (1944), and many fundamental results were established by Stephen Cole Kleene and Post (1954). The Turing degrees have been an area of intense research since then. Many proofs in the area make use of a proof technique known as the priority method.
Turing equivalence:
For the rest of this article, the word set will refer to a set of natural numbers. A set X is said to be Turing reducible to a set Y if there is an oracle Turing machine that decides membership in X when given an oracle for membership in Y. The notation X ≤T Y indicates that X is Turing reducible to Y.
Turing equivalence:
Two sets X and Y are defined to be Turing equivalent if X is Turing reducible to Y and Y is Turing reducible to X. The notation X ≡T Y indicates that X and Y are Turing equivalent. The relation ≡T can be seen to be an equivalence relation, which means that for all sets X, Y, and Z: X ≡T X X ≡T Y implies Y ≡T X If X ≡T Y and Y ≡T Z then X ≡T Z.A Turing degree is an equivalence class of the relation ≡T. The notation [X] denotes the equivalence class containing a set X. The entire collection of Turing degrees is denoted D The Turing degrees have a partial order ≤ defined so that [X] ≤ [Y] if and only if X ≤T Y. There is a unique Turing degree containing all the computable sets, and this degree is less than every other degree. It is denoted 0 (zero) because it is the least element of the poset D . (It is common to use boldface notation for Turing degrees, in order to distinguish them from sets. When no confusion can occur, such as with [X], the boldface is not necessary.) For any sets X and Y, X join Y, written X ⊕ Y, is defined to be the union of the sets {2n : n ∈ X} and {2m+1 : m ∈ Y}. The Turing degree of X ⊕ Y is the least upper bound of the degrees of X and Y. Thus D is a join-semilattice. The least upper bound of degrees a and b is denoted a ∪ b. It is known that D is not a lattice, as there are pairs of degrees with no greatest lower bound.
Turing equivalence:
For any set X the notation X′ denotes the set of indices of oracle machines that halt (when given their index as input) when using X as an oracle. The set X′ is called the Turing jump of X. The Turing jump of a degree [X] is defined to be the degree [X′]; this is a valid definition because X′ ≡T Y′ whenever X ≡T Y. A key example is 0′, the degree of the halting problem.
Basic properties of the Turing degrees:
Every Turing degree is countably infinite, that is, it contains exactly ℵ0 sets.
There are 2ℵ0 distinct Turing degrees.
For each degree a the strict inequality a < a′ holds.
For each degree a, the set of degrees below a is countable. The set of degrees greater than a has size 2ℵ0
Structure of the Turing degrees:
A great deal of research has been conducted into the structure of the Turing degrees. The following survey lists only some of the many known results. One general conclusion that can be drawn from the research is that the structure of the Turing degrees is extremely complicated.
Order properties There are minimal degrees. A degree a is minimal if a is nonzero and there is no degree between 0 and a. Thus the order relation on the degrees is not a dense order.
The Turing degrees are not linearly ordered by ≤T.
In fact, for every nonzero degree a there is a degree b incomparable with a.
There is a set of 2ℵ0 pairwise incomparable Turing degrees.
There are pairs of degrees with no greatest lower bound. Thus D is not a lattice.
Every countable partially ordered set can be embedded in the Turing degrees.
An infinite strictly increasing sequence a1, a2, ... of Turing degrees cannot have the least upper bound, but it always has an exact pair c, d such that ∀e (e<c∧e<d ⇔ ∃i e≤ai), and thus it has (non-unique) minimal upper bounds.
Assuming the axiom of constructibility, it can be shown there is a maximal chain of degrees of order type ω1 Properties involving the jump For every degree a there is a degree strictly between a and a′. In fact, there is a countable family of pairwise incomparable degrees between a and a′.
Jump inversion: a degree a is of the form b′ if and only if 0′ ≤ a.
For any degree a there is a degree b such that a < b and b′ = a′; such a degree b is called low relative to a.
There is an infinite sequence ai of degrees such that a′i+1 ≤ ai for each i.
Post's theorem establishes a close correspondence between the arithmetical hierarchy and finitely iterated Turing jumps of the empty set.
Logical properties Simpson (1977) showed that the first-order theory of D in the language ⟨ ≤, = ⟩ or ⟨ ≤, ′, = ⟩ is many-one equivalent to the theory of true second-order arithmetic. This indicates that the structure of D is extremely complicated.
Shore and Slaman (1999) showed that the jump operator is definable in the first-order structure of D with the language ⟨ ≤, = ⟩.
Recursively enumerable Turing degrees:
A degree is called recursively enumerable (r.e.) or computably enumerable (c.e.) if it contains a recursively enumerable set. Every r.e. degree is below 0′, but not every degree below 0′ is r.e.. However, a set A is many-one reducible to 0′ iff A is r.e..
(G. E. Sacks, 1964) The r.e. degrees are dense; between any two r.e. degrees there is a third r.e. degree.
(A. H. Lachlan, 1966a and C. E. M. Yates, 1966) There are two r.e. degrees with no greatest lower bound in the r.e. degrees.
(A. H. Lachlan, 1966a and C. E. M. Yates, 1966) There is a pair of nonzero r.e. degrees whose greatest lower bound is 0.
(A. H. Lachlan, 1966b) There is no pair of r.e. degrees whose greatest lower bound is 0 and whose least upper bound is 0′. This result is informally called the nondiamond theorem.
(S. K. Thomason, 1971) Every finite distributive lattice can be embedded into the r.e. degrees. In fact, the countable atomless Boolean algebra can be embedded in a manner that preserves suprema and infima.
(A. H. Lachlan and R. I. Soare, 1980) Not all finite lattices can be embedded in the r.e. degrees (via an embedding that preserves suprema and infima). A particular example is shown to the right.
Recursively enumerable Turing degrees:
(L. A. Harrington and T. A. Slaman, see Nies, Shore, and Slaman (1998)) The first-order theory of the r.e. degrees in the language ⟨ 0, ≤, = ⟩ is many-one equivalent to the theory of true first-order arithmetic.Additionally, there is Shoenfield's limit lemma, a set A satisfies [A]≤T∅′ iff there is a "recursive approximation" to its characteristic function: a function g such that for sufficiently large s, g(s)=χA(s) .A set A is called n-r e. if there is a family of functions (As)s∈N such that: As is a recursive approximation of A: for some t, for any s≥t we have As(x) = A(x), in particular conflating A with its characteristic function. (Removing this condition yields a definition of A being "weakly n-r.e.") As is an "n-trial predicate": for all x, A0(x)=0 and the cardinality of {s∣As(x)≠As+1(x)} is ≤n.Properties of n-r.e. degrees: The class of sets of n-r.e. degree is a strict subclass of the class of sets of (n+1)-r.e. degree.
Recursively enumerable Turing degrees:
For all n>1 there are two (n+1)-r.e. degrees a, b with a≤Tb , such that the segment {c∣a≤Tc≤Tb} contains no n-r.e. degrees.
A and A¯ are (n+1)-r.e. iff both sets are weakly-n-r.e.
Post's problem and the priority method:
Emil Post studied the r.e. Turing degrees and asked whether there is any r.e. degree strictly between 0 and 0′. The problem of constructing such a degree (or showing that none exist) became known as Post's problem. This problem was solved independently by Friedberg and Muchnik in the 1950s, who showed that these intermediate r.e. degrees do exist (Friedberg–Muchnik theorem). Their proofs each developed the same new method for constructing r.e. degrees, which came to be known as the priority method. The priority method is now the main technique for establishing results about r.e. sets.
Post's problem and the priority method:
The idea of the priority method for constructing a r.e. set X is to list a countable sequence of requirements that X must satisfy. For example, to construct a r.e. set X between 0 and 0′ it is enough to satisfy the requirements Ae and Be for each natural number e, where Ae requires that the oracle machine with index e does not compute 0′ from X and Be requires that the Turing machine with index e (and no oracle) does not compute X. These requirements are put into a priority ordering, which is an explicit bijection of the requirements and the natural numbers. The proof proceeds inductively with one stage for each natural number; these stages can be thought of as steps of time during which the set X is enumerated. At each stage, numbers may be put into X or forever (if not injured) prevented from entering X in an attempt to satisfy requirements (that is, force them to hold once all of X has been enumerated). Sometimes, a number can be enumerated into X to satisfy one requirement but doing this would cause a previously satisfied requirement to become unsatisfied (that is, to be injured). The priority order on requirements is used to determine which requirement to satisfy in this case. The informal idea is that if a requirement is injured then it will eventually stop being injured after all higher priority requirements have stopped being injured, although not every priority argument has this property. An argument must be made that the overall set X is r.e. and satisfies all the requirements. Priority arguments can be used to prove many facts about r.e. sets; the requirements used and the manner in which they are satisfied must be carefully chosen to produce the required result.
Post's problem and the priority method:
For example, a simple (and hence noncomputable r.e.) low X (low means X′=0′) can be constructed in infinitely many stages as follows. At the start of stage n, let Tn be the output (binary) tape, identified with the set of cell indices where we placed 1 so far (so X=∪n Tn; T0=∅); and let Pn(m) be the priority for not outputting 1 at location m; P0(m)=∞. At stage n, if possible (otherwise do nothing in the stage), pick the least i<n such that ∀m Pn(m)≠i and Turing machine i halts in <n steps on some input S⊇Tn with ∀m∈S\Tn Pn(m)≥i. Choose any such (finite) S, set Tn+1=S, and for every cell m visited by machine i on S, set Pn+1(m) = min(i, Pn(m)), and set all priorities >i to ∞, and then set one priority ∞ cell (any will do) not in S to priority i. Essentially, we make machine i halt if we can do so without upsetting priorities <i, and then set priorities to prevent machines >i from disrupting the halt; all priorities are eventually constant.
Post's problem and the priority method:
To see that X is low, machine i halts on X iff it halts in <n steps on some Tn such that machines <i that halt on X do so <n-i steps (by recursion, this is uniformly computable from 0′). X is noncomputable since otherwise a Turing machine could halt on Y iff Y\X is nonempty, contradicting the construction since X excludes some priority i cells for arbitrarily large i; and X is simple because for each i the number of priority i cells is finite. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Multiple baseline design**
Multiple baseline design:
A multiple baseline design is used in medical, psychological, and biological research. The multiple baseline design was first reported in 1960 as used in basic operant research. It was applied in the late 1960s to human experiments in response to practical and ethical issues that arose in withdrawing apparently successful treatments from human subjects. In it two or more (often three) behaviors, people or settings are plotted in a staggered graph where a change is made to one, but not the other two, and then to the second, but not the third behavior, person or setting. Differential changes that occur to each behavior, person or in each setting help to strengthen what is essentially an AB design with its problematic competing hypotheses.
Multiple baseline design:
Because treatment is started at different times, changes are attributable to the treatment rather than to a chance factor. By gathering data from many subjects (instances), inferences can be made about the likeliness that the measured trait generalizes to a greater population. In multiple baseline designs, the experimenter starts by measuring a trait of interest, then applies a treatment before measuring that trait again. Treatment does not begin until a stable baseline has been recorded, and does not finish until measures regain stability. If a significant change occurs across all participants the experimenter may infer that the treatment is effective.
Multiple baseline design:
Multiple base-line experiments are most commonly used in cases where the dependent variable is not expected to return to normal after the treatment has been applied, or when medical reasons forbid the withdrawal of a treatment. They often employ particular methods or recruiting participants. Multiple baseline designs are associated with potential confounds introduced by experimenter bias, which must be addressed to preserve objectivity. Particularly, researchers are advised to develop all test schedules and data collection limits beforehand.
Recruiting participants:
Although multiple baseline designs may employ any method of recruitment, it is often associated with "ex post facto" recruitment. This is because multiple baselines can provide data regarding the consensus of a treatment response. Such data can often not be gathered from ABA (reversal) designs for ethical or learning reasons. Experimenters are advised not to remove cases that do not exactly fit their criteria, as this may introduce sampling bias and threaten validity.
Recruiting participants:
Ex post facto recruitment methods are not considered true experiments, due to the limits of experimental control or randomized control that the experimenter has over the trait. This is because a control group may necessarily be selected from a discrete separate population. This research design is thus considered a quasi-experimental design.
Concurrent designs:
Multiple baseline studies are often categorized as either concurrent or nonconcurrent. Concurrent designs are the traditional approach to multiple baseline studies, where all participants undergo treatment simultaneously. This strategy is advantageous because it moderates several threats to validity, and history effects in particular. Concurrent multiple baseline designs are also useful for saving time, since all participants are processed at once. The ability to retrieve complete data sets within well defined time constraints is a valuable asset while planning research.
Nonconcurrent designs:
Nonconcurrent multiple baseline studies apply treatment to several individuals at delayed intervals. This has the advantage of greater flexibility in recruitment of participants and testing location. For this reason, perhaps, nonconcurrent multiple baseline experiments are recommended for research in an educational setting.
It is recommended that the experimenter selects time frames beforehand to avoid experimenter bias, but even when methods are used to improve validity, inferences may be weakened. Currently, there is debate as to whether nonconcurrent studies represent a real threat from history effects. It is generally agreed, however, that concurrent testing is more stable.
Disadvantages:
Although multiple baseline experimental designs compensate for many of the issues inherent in ex post facto recruitment, experimental manipulation of a trait gathered by this method may not be manipulated. Thus these studies are prevented from inferring causation if there are no phases to demonstrate reversibility. However, if such phases are included (as is the standard of experimentation), they can successfully demonstrate causation.
Managing threats to validity:
A priori (beforehand) specification of the hypothesis, time frames, and data limits help control threats due to experimenter bias. For the same reason researchers should avoid removing participants based on merit. Multiple probe designs may be useful in identifying extraneous factors which may be influencing your results. Lastly, experimenters should avoid gathering data during sessions alone. If in-session data is gathered a note of the dates should be tagged to each measurement in order to provide an accurate time-line for potential reviewers. This data may represent unnatural behaviour or states of mind, and must be considered carefully during interpretation. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Statistics education**
Statistics education:
Statistics education is the practice of teaching and learning of statistics, along with the associated scholarly research.
Statistics education:
Statistics is both a formal science and a practical theory of scientific inquiry, and both aspects are considered in statistics education. Education in statistics has similar concerns as does education in other mathematical sciences, like logic, mathematics, and computer science. At the same time, statistics is concerned with evidence-based reasoning, particularly with the analysis of data. Therefore, education in statistics has strong similarities to education in empirical disciplines like psychology and chemistry, in which education is closely tied to "hands-on" experimentation.
Statistics education:
Mathematicians and statisticians often work in a department of mathematical sciences (particularly at colleges and small universities). Statistics courses have been sometimes taught by non-statisticians, against the recommendations of some professional organizations of statisticians and of mathematicians.
Statistics education research is an emerging field that grew out of different disciplines and is currently establishing itself as a unique field that is devoted to the improvement of teaching and learning statistics at all educational levels.
Goals of statistics education:
Statistics educators have cognitive and noncognitive goals for students. For example, former American Statistical Association (ASA) President Katherine Wallman defined statistical literacy as including the cognitive abilities of understanding and critically evaluating statistical results as well as appreciating the contributions statistical thinking can make.
Goals of statistics education:
Cognitive goals In the text rising from the 2008 joint conference of the International Commission on Mathematical Instruction and the International Association of Statistics Educators, editors Carmen Batanero, Gail Burrill, and Chris Reading (Universidad de Granada, Spain, Michigan State University, USA, and University of New England, Australia, respectively) note worldwide trends in curricula which reflect data-oriented goals. In particular, educators currently seek to have students: "design investigations; formulate research questions; collect data using observations, surveys, and experiments; describe and compare data sets; and propose and justify conclusions and predictions based on data." The authors note the importance of developing statistical thinking and reasoning in addition to statistical knowledge.
Goals of statistics education:
Despite the fact that cognitive goals for statistics education increasingly focus on statistical literacy, statistical reasoning, and statistical thinking rather than on skills, computations and procedures alone, there is no agreement about what these terms mean or how to assess these outcomes. A first attempt to define and distinguish between these three terms appears in the ARTIST website which was created by Garfield, delMas and Chance and has since been included in several publications.
Goals of statistics education:
Brief definitions of these terms are as follows: Statistical literacy is being able to read and use basic statistical language and graphical representations to understand statistical information in the media and in daily life.
Statistical reasoning is being able to reason about and connect different statistical concepts and ideas, such as knowing how and why outliers affect statistical measures of center and variability.
Goals of statistics education:
Statistical thinking is the type of thinking used by statisticians when they encounter a statistical problem. This involves thinking about the nature and quality of the data and, where the data came from, choosing appropriate analyses and models, and interpreting the results in the context of the problem and given the constraints of the data.Further cognitive goals of statistics education vary across students' educational level and the contexts in which they expect to encounter statistics.
Goals of statistics education:
Statisticians have proposed what they consider the most important statistical concepts for educated citizens. For example, Utts (2003) published seven areas of what every educated citizen should know, including understanding that "variability is normal" and how "coincidences… are not uncommon because there are so many possibilities." Gal (2002) suggests adults in industrialized societies are expected to exercise statistical literacy, "the ability to interpret and critically evaluate statistical information… in diverse contexts, and the ability to… communicate understandings and concerns regarding the… conclusions." Non-cognitive goals Non-cognitive outcomes include affective constructs such as attitudes, beliefs, emotions, dispositions, and motivation. According to prominent researchers Gal & Ginsburg, statistics educators should make it a priority to be aware of students' ideas, reactions, and feelings towards statistics and how these affect their learning.
Goals of statistics education:
Beliefs Beliefs are defined as one's individually held ideas about statistics, about oneself as a learner of statistics, and about the social context of learning statistics. Beliefs are distinct from attitudes in the sense that attitudes are relatively stable and intense feelings that develop over time in the context of experiences learning statistics. Students' web of beliefs provides a context for their approach towards their classroom experiences in statistics. Many students enter a statistics course with apprehension towards learning the subject, which works against the learning environment that the instructor is trying to accomplish. Therefore, it is important for instructors to have access to assessment instruments that can give an initial diagnosis of student beliefs and monitor beliefs during a course. Frequently, assessment instruments have monitored beliefs and attitudes together. For examples of such instruments, see the attitudes section below.
Goals of statistics education:
Dispositions Disposition has to do with the ways students question the data and approach a statistical problem. Dispositions is one of the four dimensions in Wild and Pfannkuch's framework for statistical thinking, and contains the following elements: Curiosity and Awareness: These traits are a part of the process of generating questions and generating ideas to explore and analyze data.
Engagement: Students will be most observant and aware in the areas they find most interesting.
Imagination: This trait is important for viewing a problem from different perspectives and coming up with possible explanations.
Scepticism: Critical thinking is important for receiving new ideas and information and evaluating the appropriateness of study design and analysis.
Being logical: The ability to detect when one idea follows from another is important for arriving at valid conclusions.
Goals of statistics education:
A propensity to seek deeper meaning: This means not taking everything at face value and being open to consider new ideas and dig deeper for information.Scheaffer states that a goal of statistics education is to have students see statistics broadly. He developed a list of views of statistics that can lead to this broad view, and describes them as follows: Statistics as number sense: Do I understand what the numbers mean? (seeing data as numbers in context, reading charts, graphs and tables, understanding numerical and graphical summaries of data, etc.) Statistics as a way of understanding the world: Can I use existing data to help make decisions? (using census data, birth and death rates, disease rates, CPI, ratings, rankings, etc., to describe, decide and defend) Statistics as organized problem solving: Can I design and carry out a study to answer specific questions? (pose problem, collect data according to a plan, analyze data, and draw conclusions from data) Attitudes Since students often experience math anxiety and negative opinions about statistics courses, various researchers have addressed attitudes and anxiety towards statistics. Some instruments have been developed to measure college students' attitudes towards statistics, and have been shown to have appropriate psychometric properties. Examples of such instruments include: Survey of Attitudes Towards Statistics (SATS), developed by Schau, Stevens, Dauphinee, and Del Vecchio Attitude Toward Statistics Scale, developed by Wise Statistics Attitude Survey (SAS), developed by Roberts and BilderbackCareful use of instruments such as these can help statistics instructors to learn about students' perception of statistics, including their anxiety towards learning statistics, the perceived difficulty of learning statistics, and their perceived usefulness of the subject. Some studies have shown modest success at improving student attitudes in individual courses, but no generalizable studies showing improvement in student attitudes have been seen.
Goals of statistics education:
Nevertheless, one of the goals of statistics education is to make the study of statistics a positive experience for students and to bring in interesting and engaging examples and data that will motivate students. According to a fairly recent literature review, improved student attitudes towards statistics can lead to better motivation and engagement, which also improves cognitive learning outcomes.
Primary–secondary education level:
New Zealand In New Zealand, a new curriculum for statistics has been developed by Chris Wild and colleagues at Auckland University. Rejecting the contrived, and now unnecessary due to computer power, approach of reasoning under the null and the restrictions of normal theory, they use comparative box plots and bootstrap to introduce concepts of sampling variability and inference. The developing curriculum also contains aspects of statistical literacy.
Primary–secondary education level:
United Kingdom In the United Kingdom, at least some statistics has been taught in schools since the 1930s. At present, A-level qualifications (typically taken by 17- to 18-year-olds) are being developed in "Statistics" and "Further Statistics". The coverage of the former includes: Probability; Data Collection; Descriptive Statistics; Discrete Probability Distributions; Binomial Distribution; Poisson Distributions; Continuous Probability Distributions; The Normal Distribution; Estimation; Hypothesis Testing; Chi-Squared; Correlation and Regression. The coverage of "Further Statistics" includes: Continuous Probability Distributions; Estimation; Hypothesis Testing; One Sample Tests; Hypothesis Testing; Two Sample Tests; Goodness of Fit Tests; Experimental Design; Analysis of Variance (Anova); Statistical Process Control; Acceptance Sampling. The Centre for Innovation in Mathematics Teaching (CIMT) has online course notes for these sets of topics. Revision notes for an existing qualification indicate a similar coverage. At an earlier age (typically 15–16 years) GCSE qualifications in mathematics contain "Statistics and Probability" topics on: Probability; Averages; Standard Deviation; Sampling; Cumumulative Frequency Graphs (including median and quantiles); Representing Data; Histograms. The UK's Office for National Statistics has a webpage leading to material suitable for both teachers and students at school level. In 2004 the Smith inquiry made the following statement: "There is much concern and debate about the positioning of Statistics and Data Handling within the current mathematics GCSE, where it occupies some 25 per cent of the timetable allocation. On the one hand, there is widespread agreement that the Key Stage 4 curriculum is over-crowded and that the introduction of Statistics and Data Handling may have been at the expense of time needed for practising and acquiring fluency in core mathematical manipulations. Many in higher education mathematics and engineering departments take this view. On the other hand, there is overwhelming recognition, shared by the Inquiry, of the vital importance of Statistics and Data Handling skills both for a number of other academic disciplines and in the workplace. The Inquiry recommends that there be a radical re-look at this issue and that much of the teaching and learning of Statistics and Data Handling would be better removed from the mathematics timetable and integrated with the teaching and learning of other disciplines (e.g. biology or geography). The time restored to the mathematics timetable should be used for acquiring greater mastery of core mathematical concepts and operations." United States In the United States, schooling has increased the use of probability and statistics, especially since the 1990s. Summary statistics and graphs are taught in elementary school in many states. Topics in probability and statistical reasoning are taught in high school algebra (or mathematical science) courses; statistical reasoning has been examined in the SAT test since 1994. The College Board has developed an Advanced Placement course in statistics, which has provided a college-level course in statistics to hundreds of thousands of high school students, with the first examination happening in May 1997. In 2007, the ASA endorsed the Guidelines for Assessment and Instruction in Statistics Education (GAISE), a two-dimensional framework for the conceptual understanding of statistics in Pre-K-12 students. The framework contains learning objectives for students at each conceptual level and provides pedagogical examples that are consistent with the conceptual levels.
Primary–secondary education level:
Estonia Estonia is piloting a new statistics curriculum developed by the Computer-Based Math foundation based around its principles of using computers as the primary tool of education. in cooperation with the University of Tartu.
University level:
General Statistics is often taught in departments of mathematics or in departments of mathematical sciences. At the undergraduate level, statistics is often taught as a service course.
University level:
United Kingdom By tradition in the U.K., most professional statisticians are trained at the Master level. A difficulty of recruiting strong undergraduates has been noted: "Very few undergraduates positively choose to study statistics degrees; most choose some statistics options within a mathematics programme, often to avoid the advanced pure and applied mathematics courses. My view is that statistics as a theoretical discipline is better taught late rather than early, whereas statistics as part of scientific methodology should be taught as part of science."In the United Kingdom, the teaching of statistics at university level was originally done within science departments that needed the topic to accompany the teaching of their own subjects, and departments of mathematics had limited coverage before the 1930s. For the twenty years subsequent to this, while departments of mathematics had started to teach statistics, there was little realisation that essentially the same basic statistical methodology was being applied across a variety of sciences. Statistical departments have had difficulty when they have been separated from mathematics departments.Psychologist Andy Field (British Psychological Society Teaching and Book Award) created a new concept of statistical teaching and textbooks that goes beyond the printed page.
University level:
United States Enrollments in statistics have increased in community colleges, in four-year colleges and universities in the United States. At community colleges in the United States, mathematics has experienced increased enrollment since 1990. At community colleges, the ratio of the students enrolled in statistics to those enrolled in calculus rose from 56% in 1990 to 82% in 1995. One of the ASA-endorsed GAISE reports focused on statistics education at the introductory college level. The report includes a brief history of the introductory statistics course and recommendations for how it should be taught.
University level:
In many colleges, a basic course in "statistics for non-statisticians" has required only algebra (and not calculus); for future statisticians, in contrast, the undergraduate exposure to statistics is highly mathematical. As undergraduates, future statisticians should have completed courses in multivariate calculus, linear algebra, computer programming, and a year of calculus-based probability and statistics. Students wanting to obtain a doctorate in statistics from "any of the better graduate programs in statistics" should also take "real analysis". Laboratory courses in physics, chemistry and psychology also provide useful experiences with planning and conducting experiments and with analyzing data. The ASA recommends that undergraduate students consider obtaining a bachelor's degree in applied mathematics as preparation for entering a master program in statistics.Historically, professional degrees in statistics have been at the Master level, although some students may qualify to work with a bachelor's degree and job-related experience or further self-study. Professional competence requires a background in mathematics—including at least multivariate calculus, linear algebra, and a year of calculus-based probability and statistics. In the United States, a master program in statistics requires courses in probability, mathematical statistics, and applied statistics (e.g., design of experiments, survey sampling, etc.).
University level:
For a doctoral degree in statistics, it has been traditional that students complete a course in measure-theoretic probability as well as courses in mathematical statistics. Such courses require a good course in real analysis, covering the proofs of the theory of calculus and topics like the uniform convergence of functions. In recent decades, some departments have discussed allowing doctoral students to waive the course in measure-theoretic probability by demonstrating advanced skills in computer programming or scientific computing.
Who should teach statistics?:
The question of what qualities are needed to teach statistics has been much discussed, and sometimes this discussion is concentrated on the qualifications necessary for those undertaking such teaching. The question arises separately for teaching at both school and university levels, partly because of the need for numerically more such teachers at school level and partly because of need for such teachers to cover a broad range of other topics within their overall duties. Given that "statistics" is often taught to non-scientists, opinions can range all the way from "statistics should be taught by statisticians", through "teaching of statistics is too mathematical" to the extreme that "statistics should not be taught by statisticians".
Who should teach statistics?:
Teaching at university level In the United States especially, statisticians have long complained that many mathematics departments have assigned mathematicians (without statistical competence) to teach statistics courses, effectively giving "double blind" courses. The principle that college-instructors should have qualifications and engagement with their academic discipline has long been violated in United States colleges and universities, according to generations of statisticians. For example, the journal Statistical Science reprinted "classic" articles on the teaching of statistics by non-statisticians by Harold Hotelling; Hotelling's articles are followed by the comments of Kenneth J. Arrow, W. Edwards Deming, Ingram Olkin, David S. Moore, James V. Sidek, Shanti S. Gupta, Robert V. Hogg, Ralph A. Bradley, and by Harold Hotelling, Jr. (an economist and son of Harold Hotelling).
Who should teach statistics?:
Data on the teaching of statistics in the United States has been collected on behalf of the Conference Board of the Mathematical Sciences (CBMS). Examining data from 2000, Schaeffer and Stasny reported By far the majority of instructors within statistics departments have at least a master’s degree in statistics or biostatistics (about 89% for doctoral departments and about 79% for master’s departments). In doctoral mathematics departments, however, only about 58% of statistics course instructors had at least a master’s degree in statistics or biostatistics as their highest degree earned. In master’s-level mathematics departments, the corresponding percentage was near 44%, and in bachelor’s-level departments only 19% of statistics course instructors had at least a master’s degree in statistics or biostatistics as their highest degree earned. As we expected, a large majority of instructors in statistics departments (83% for doctoral departments and 62% for master’s departments) held doctoral degrees in either statistics or biostatistics. The comparable percentages for instructors of statistics in mathematics departments were about 52% and 38%.
Who should teach statistics?:
The principle that statistics-instructors should have statistical competence has been affirmed by the guidelines of the Mathematical Association of America, which has been endorsed by the ASA. The unprofessional teaching of statistics by mathematicians (without qualifications in statistics) has been addressed in many articles.
Teaching methods:
The literature on methods of teaching statistics is closely related to the literature on the teaching of mathematics for two reasons. First, statistics is often taught as part of the mathematics curriculum, by instructors trained in mathematics and working in a mathematics department. Second, statistical theory has often been taught as a mathematical theory rather than as the practical logic of science --- as the science that "puts chance to work" in Rao's phrase--- and this has entailed an emphasis on formal and manipulative training, such as solving combinatorial problems involving red and green jelly beans. Statisticians have complained that mathematicians are prone to over-emphasize mathematical manipulations and probability theory and under-emphasize questions of experimentation, survey methodology, exploratory data analysis, and statistical inference.In recent decades, there has been an increased emphasis on data analysis and scientific inquiry in statistics education. In the United Kingdom, the Smith inquiry Making Mathematics Count suggests teaching basic statistical concepts as part of the science curriculum, rather than as part of mathematics. In the United States, the ASA's guidelines for undergraduate statistics specify that introductory statistics should emphasize the scientific methods of data collection, particularly randomized experiments and random samples: further, the first course should review these topics when the theory of "statistical inference" is studied. Similar recommendations occur for the Advanced Placement (AP) course in Statistics. The ASA and AP guidelines are followed by contemporary textbooks in the US, such as those by Freedman, Purvis & Pisani (Statistics) and by David S. Moore (Introduction to the Practice of Statistics with McCabe and Statistics: Concepts and Controversies with Notz) and by Watkins, Schaeffer & Cobb (Statistics: From Data to Decisions and Statistics in Action).
Teaching methods:
Besides an emphasis on the scientific inquiry in the content of beginning of statistics, there has also been an increase on active learning in the conduct of the statistics classroom.
Professional community:
Associations The International Statistical Institute (ISI) now has one section devoted to education, the International Association for Statistical Education (IASE), which runs the International Conference on Teaching Statistics every four years as well as IASE satellite conferences around ISI and ICMI meetings. The UK established the Royal Statistical Society Centre for Statistics Education and the ASA now also has a Section on Statistical Education, focused mostly on statistics teaching at the elementary and secondary levels.
Professional community:
Conferences In addition to the international gatherings of statistics educators at ICOTS every four years, the US hosts a US Conference on Teaching Statistics (USCOTS) every two years and has recently started an Electronic Conference on Teaching Statistics (eCOTS) to alternate with USCOTS. Sessions on statistics education area also offered at many conferences in mathematics educations such as the International Congress on Mathematical Education, the National Council of Teachers of Mathematics, the Conference of the International Group for the Psychology of Mathematics Education, and the Mathematics Education Research Group of Australasia. The annual Joint Statistical Meetings (offered by the ASA and Statistics Canada) offer many sessions and roundtables on statistics education. The International Research Forums on Statistical Reasoning, Thinking, and Literacy offer scientific gatherings every two years and related publications in journals, CD-ROMs and books on research in statistics education.
Professional community:
Graduate coursework and programs Only three universities currently offer graduate programs in statistics education: the University of Granada, the University of Minnesota, and the University of Florida. However, graduate students in a variety of disciplines (e.g., mathematics education, psychology, educational psychology) have been finding ways to complete dissertations on topics related to teaching and learning statistics. These dissertations are archived on the IASE web site.Two main courses in statistics education that have been taught in a variety of settings and departments are a course on teaching statistics and a course on statistics education research. An ASA-sponsored workshop has established recommendations for additional graduate programs and courses.
Software for learning:
Fathom: Dynamic Data Software TinkerPlots StatCrunch
Trends in Statistics Education:
Teachers of statistics have been encouraged to explore new directions in curriculum content, pedagogy and assessment. In an influential talk at USCOTS, researcher George Cobb presented an innovative approach to teaching statistics that put simulation, randomization, and bootstrapping techniques at the core of the college-level introductory course, in place of traditional content such as probability theory and the t-test. Several teachers and curriculum developers have been exploring ways to introduce simulation, randomization, and bootstrapping as teaching tools for the secondary and postsecondary levels. Courses such as the University of Minnesota's CATALST, Nathan Tintle and collaborators' Introduction to Statistical Investigations, and the Lock team's Unlocking the Power of Data, are curriculum projects based on Cobb's ideas. Other researchers have been exploring the development of informal inferential reasoning as a way to use these methods to build a better understanding of statistical inference.
Trends in Statistics Education:
Another recent direction is addressing the big data sets that are increasingly affecting or being contributed to in our daily lives. Statistician Rob Gould, creator of Data Cycle, The Musical dinner and theatre spectacular, outlines many of these types of data and encourages teachers to find ways to use the data and address issues around big data. According to Gould, curricula focused on big data will address issues of sampling, prediction, visualization, data cleaning, and the underlying processes that generate data, rather than traditionally emphasized methods of making statistical inferences such as hypothesis testing.
Trends in Statistics Education:
Driving both of these changes is the increased role of computing in teaching and learning statistics. Some researchers argue that as the use of modeling and simulation increase, and as data sets become larger and more complex, students will need better and more technical computing skills. Projects such as MOSAIC have been creating courses that blend computer science, modeling, and statistics. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Stimulant psychosis**
Stimulant psychosis:
Stimulant psychosis is a mental disorder characterized by psychotic symptoms (such as hallucinations, paranoid ideation, delusions, disorganized thinking, grossly disorganized behaviour). It involves and typically occurs following an overdose or several day 'binge' on psychostimulants; however, one study reported occurrences at regularly prescribed doses in approximately 0.1% of individuals within the first several weeks after starting amphetamine or methylphenidate therapy. Methamphetamine psychosis, or long-term effects of stimulant use in the brain (at the molecular level), depend upon genetics and may persist for some time.The most common causative agents are substituted amphetamines, including substituted cathinones, as well as certain dopamine reuptake inhibitors such as cocaine and phenidates.
Signs and symptoms:
The symptoms of stimulant psychosis vary depending on the drug ingested, but generally involve the symptoms of organic psychosis such as hallucinations, delusions, or paranoia. Other symptoms may include mania, erratic behavior, agitation and/or aggression.
Cause:
Substituted amphetamines Drugs in the class of amphetamines, or substituted amphetamines, are known to induce "amphetamine psychosis" typically when chronically abused or used in high doses. In an Australian study of 309 active methamphetamine users, 18% had experienced a clinical level psychosis in the past year. Commonly abused amphetamines include methamphetamine, MDMA, 4-FA, as well as substituted cathinones like a-PVP, MDPV, and mephedrone, though a large number of other closely related compounds have been recently synthesized. Methylphenidate is sometimes incorrectly included in this class, although it is nonetheless still capable of producing stimulant psychosis.
Cause:
The symptoms of amphetamine psychosis include auditory and visual hallucinations, grandiosity, delusions of persecution, and delusions of reference concurrent with both clear consciousness and prominent extreme agitation. A Japanese study of recovery from methamphetamine psychosis reported a 64% recovery rate within 10 days rising to an 82% recovery rate at 30 days after methamphetamine cessation. However it has been suggested that around 5–15% of users fail to make a complete recovery in the long term. Furthermore, even at a small dose, the psychosis can be quickly reestablished. Psychosocial stress has been found to be an independent risk factor for psychosis relapse even without further substituted amphetamine use in certain cases.The symptoms of acute amphetamine psychosis are very similar to those of the acute phase of schizophrenia although in amphetamine psychosis visual hallucinations are more common and thought disorder is rare. Amphetamine psychosis may be purely related to high drug usage, or high drug usage may trigger an underlying vulnerability to schizophrenia. There is some evidence that vulnerability to amphetamine psychosis and schizophrenia may be genetically related. Relatives of methamphetamine users with a history of amphetamine psychosis are five times more likely to have been diagnosed with schizophrenia than relatives of methamphetamine users without a history of amphetamine psychosis. The disorders are often distinguished by a rapid resolution of symptoms in amphetamine psychosis, while schizophrenia is more likely to follow a chronic course.Although rare and not formally recognized, a condition known as Amphetamine Withdrawal Psychosis (AWP) may occur upon cessation of substituted amphetamine use and, as the name implies, involves psychosis that appears on withdrawal from substituted amphetamines. However, unlike similar disorders, in AWP, substituted amphetamines reduce rather than increase symptoms, and the psychosis or mania resolves with resumption of the previous dosing schedule.
Cause:
Cocaine Cocaine has a similar potential to induce temporary psychosis with more than half of cocaine abusers reporting at least some psychotic symptoms at some point. Typical symptoms include paranoid delusions that they are being followed and that their drug use is being watched, accompanied by hallucinations that support the delusional beliefs. Delusional parasitosis with formication ("cocaine bugs") is also a fairly common symptom.Cocaine-induced psychosis shows sensitization toward the psychotic effects of the drug. This means that psychosis becomes more severe with repeated intermittent use.
Cause:
Phenidates Methylphenidate and its analogues (such as ethylphenidate, 4F-MPH, and isopropylphenidate) share similar pharmacological profiles as other norepinephrine-dopamine reuptake inhibitors. Chronic abuse of methylphenidate has the potential to lead to psychosis. Similar psychiatric side effects have been reported in a study of ethylphenidate. No studies regarding psychosis and 4F-MPH or isopropylphenidate have been conducted, but given their high DAT binding and cellular uptake activity, the possibility of stimulant psychosis remains.
Cause:
Caffeine There is limited evidence that caffeine, in high doses or when chronically abused, may induce psychosis in normal individuals and worsen pre-existing psychosis in those diagnosed with schizophrenia.
Diagnosis:
Differential diagnosis Though less common than stimulant psychosis, stimulants such as cocaine and amphetamines as well as the dissociative drug phencyclidine (PCP, angel dust) may also cause a theorized severe and life-threatening condition known as excited delirium. This condition manifests as a combination of delirium, psychomotor agitation, anxiety, delusions, hallucinations, speech disturbances, disorientation, violent and bizarre behavior, insensitivity to pain, elevated body temperature, and hysterical strength. Despite some superficial similarities in presentation excited delirium is a distinct (and more serious) condition than stimulant psychosis. The existence of excited delirium is currently debated.
Transition to schizophrenia:
A 2019 systematic review and meta-analysis by Murrie et al. found that the pooled proportion of transition from amphetamine-induced psychosis to schizophrenia was 22% (5 studies, CI 14%–34%). This was lower than cannabis (34%) and hallucinogens (26%), but higher than opioid (12%), alcohol (10%) and sedative (9%) induced psychoses. Transition rates were slightly lower in older cohorts but were not affected by sex, country of the study, hospital or community location, urban or rural setting, diagnostic methods, or duration of follow-up.
Treatment:
Treatment consists of supportive care during the acute intoxication phase: maintaining hydration, body temperature, blood pressure, and heart rate at acceptable levels until the drug is sufficiently metabolized to allow vital signs to return to baseline. Typical and atypical antipsychotics have been shown to be helpful in the early stages of treatment. However, the benzodiazepines, temazepam and triazolam at 30 mg and 0.5 mg, respectively, are highly effective if aggression, agitation, or violent behaviour is apparent. In the instance of persistent psychosis after repeated use of stimulants, there are cases in which electroconvulsive therapy has been beneficial. This is followed by abstinence from psychostimulants supported with counselling or medication designed to assist the individual preventing a relapse and the resumption of a psychotic state. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Service provider interface**
Service provider interface:
Service provider interface (SPI) is an API intended to be implemented or extended by a third party. It can be used to enable framework extension and replaceable components.
Details:
From Java documentation: A service is a well-known set of interfaces and (usually abstract) classes. A service provider is a specific implementation of a service. The classes in a provider typically implement the interfaces and subclass the classes defined in the service itself. Service providers can be installed in an implementation of the Java platform in the form of extensions, that is, jar files placed into any of the usual extension directories. Providers can also be made available by adding them to the application's class path or by some other platform-specific means.
Details:
The concept can be extended to other platforms using the corresponding tools. In the Java Runtime Environment, SPIs are used in: Java Database Connectivity Java Cryptography Extension Java Naming and Directory Interface Java API for XML Processing Java Business Integration Java Sound Java Image I/O Java File Systems | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Contiguity (probability theory)**
Contiguity (probability theory):
In probability theory, two sequences of probability measures are said to be contiguous if asymptotically they share the same support. Thus the notion of contiguity extends the concept of absolute continuity to the sequences of measures.
The concept was originally introduced by Le Cam (1960) as part of his foundational contribution to the development of asymptotic theory in mathematical statistics. He is best known for the general concepts of local asymptotic normality and contiguity.
Definition:
Let (Ωn,Fn) be a sequence of measurable spaces, each equipped with two measures Pn and Qn. We say that Qn is contiguous with respect to Pn (denoted Qn ◁ Pn) if for every sequence An of measurable sets, Pn(An) → 0 implies Qn(An) → 0.
Definition:
The sequences Pn and Qn are said to be mutually contiguous or bi-contiguous (denoted Qn ◁▷ Pn) if both Qn is contiguous with respect to Pn and Pn is contiguous with respect to Qn.The notion of contiguity is closely related to that of absolute continuity. We say that a measure Q is absolutely continuous with respect to P (denoted Q ≪ P) if for any measurable set A, P(A) = 0 implies Q(A) = 0. That is, Q is absolutely continuous with respect to P if the support of Q is a subset of the support of P, except in cases where this is false, including, e.g., a measure that concentrates on an open set, because its support is a closed set and it assigns measure zero to the boundary, and so another measure may concentrate on the boundary and thus have support contained within the support of the first measure, but they will be mutually singular. In summary, this previous sentence's statement of absolute continuity is false. The contiguity property replaces this requirement with an asymptotic one: Qn is contiguous with respect to Pn if the "limiting support" of Qn is a subset of the limiting support of Pn. By the aforementioned logic, this statement is also false.
Definition:
It is possible however that each of the measures Qn be absolutely continuous with respect to Pn, while the sequence Qn not being contiguous with respect to Pn.
Definition:
The fundamental Radon–Nikodym theorem for absolutely continuous measures states that if Q is absolutely continuous with respect to P, then Q has density with respect to P, denoted as ƒ = dQ⁄dP, such that for any measurable set A Q(A)=∫AfdP, which is interpreted as being able to "reconstruct" the measure Q from knowing the measure P and the derivative ƒ. A similar result exists for contiguous sequences of measures, and is given by the Le Cam's third lemma.
Properties:
For the case (Pn,Qn)=(P,Q) for all n it applies Qn◃Pn⇔Q≪P It is possible that Pn≪Qn is true for all n without Pn◃Qn
Le Cam's first lemma:
For two sequences of measures and (Qn) on measurable spaces (Ωn,Fn) the following statements are equivalent: Pn◃Qn along a subsequence ⇒P(U>0)=1 along a subsequence ⇒E(V)=1 Tn⟶Pn0⇒Tn⟶Qn0 for any statistics Tn:Ωn→R .where U and V are random variables on and (Ω′,F′,Q)
Additional literature:
Roussas, George G. (1972), Contiguity of Probability Measures: Some Applications in Statistics, CUP, ISBN 978-0-521-09095-7.
Scott, D.J. (1982) Contiguity of Probability Measures, Australian & New Zealand Journal of Statistics, 24 (1), 80–88. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Angioplasty**
Angioplasty:
Angioplasty, also known as balloon angioplasty and percutaneous transluminal angioplasty (PTA), is a minimally invasive endovascular procedure used to widen narrowed or obstructed arteries or veins, typically to treat arterial atherosclerosis. A deflated balloon attached to a catheter (a balloon catheter) is passed over a guide-wire into the narrowed vessel and then inflated to a fixed size. The balloon forces expansion of the blood vessel and the surrounding muscular wall, allowing an improved blood flow. A stent may be inserted at the time of ballooning to ensure the vessel remains open, and the balloon is then deflated and withdrawn. Angioplasty has come to include all manner of vascular interventions that are typically performed percutaneously.
Angioplasty:
The word is composed of the combining forms of the Greek words ἀγγεῖον angeîon "vessel" or "cavity" (of the human body) and πλάσσω plássō "form" or "mould".
Uses and indications:
Coronary angioplasty A coronary angioplasty is a therapeutic procedure to treat the stenotic (narrowed) coronary arteries of the heart found in coronary heart disease. These stenotic segments of the coronary arteries arise due to the buildup of cholesterol-laden plaques that form in a condition known as atherosclerosis. A percutaneous coronary intervention (PCI), or coronary angioplasty with stenting, is a non-surgical procedure used to improve the blood flow to the heart.Coronary angioplasty is indicated for coronary artery disease such as unstable angina, NSTEMI, STEMI and spontaneous coronary artery perforation. PCI for stable coronary disease has been shown to significantly relieve symptoms such as angina, or chest pain, thereby improving functional limitations and quality of life.
Uses and indications:
Peripheral angioplasty Peripheral angioplasty refers to the use of a balloon to open a blood vessel outside the coronary arteries. It is most commonly done to treat atherosclerotic narrowings of the abdomen, leg and renal arteries caused by peripheral artery disease. Often, peripheral angioplasty is used in conjunction with guide wire, peripheral stenting and an atherectomy.
Uses and indications:
Chronic limb-threatening ischemia Angioplasty can be used to treat advanced peripheral artery disease to relieve the claudication, or leg pain, that is classically associated with the condition.The bypass versus angioplasty in severe ischemia of the leg (BASIL) trial investigated infrainguinal bypass surgery first compared to angioplasty first in select patients with severe lower limb ischemia who were candidates for either procedure. The BASIL trial found that angioplasty was associated with less short term morbidity compared with bypass surgery, however long term outcomes favor bypass surgery.Based on the BASIL trial, the ACCF/AHA guidelines recommend balloon angioplasty only for patients with a life expectancy of 2 years or less or those who do not have an autogenous vein available. For patients with a life expectancy greater than 2 of years life, or who have an autogenous vein, a bypass surgery could be performed first.
Uses and indications:
Renal artery angioplasty Renal artery stenosis is associated with hypertension and loss of renal function. Atherosclerotic obstruction of the renal artery can be treated with angioplasty with or without stenting of the renal artery. There is a weak recommendation for renal artery angioplasty in patients with renal artery stenosis and flash edema or congestive heart failure.
Carotid angioplasty Carotid artery stenosis can be treated with angioplasty and carotid stenting for patients at high risk for undergoing carotid endarterectomy (CEA). Although carotid endarterectomy is typically preferred over carotid artery stenting, stenting is indicated in select patients with radiation-induced stenosis or a carotid lesion not suitable for surgery.
Uses and indications:
Venous angioplasty Angioplasty is used to treat venous stenosis affecting dialysis access, with drug-coated balloon angioplasty proving to have better 6 month and 12 month patency than conventional balloon angioplasty. Angioplasty is occasionally used to treat residual subclavian vein stenosis following decompression surgery for thoracic outlet syndrome. There is a weak recommendation for deep venous stenting to treat obstructive chronic venous disease.
Contraindications:
Angioplasty requires an access vessel, typically the femoral or radial artery or femoral vein, to permit access to the vascular system for the wires and catheters used. If no access vessel of sufficient size and quality is available, angioplasty is contraindicated. A small vessel diameter, the presence of posterior calcification, occlusion, hematoma, or an earlier placement of a bypass origin, may make access to the vascular system too difficult.Percutaneous transluminal coronary angioplasty (PTCA) is contraindicated in patients with left main coronary artery disease, due to the risk of spasm of the left main coronary artery during the procedure. Also, PTCA is not recommended if there is less than 70% stenosis of the coronary arteries, as the stenosis it is not deemed to be hemodynamically significant below this level.
Technique:
Access to the vascular system is typically gained percutaneously (through the skin, without a large surgical incision). An introducer sheath is inserted into the blood vessel via the Seldinger technique. Fluoroscopic guidance uses magnetic resonance or X-ray fluoroscopy and radiopaque contrast dye to guide angled wires and catheters to the region of the body to be treated in real time. Tapered guidewire is chosen for small occlusion, followed by intermediate type guidewires for tortuous arteries and difficulty passing through extremely narrow channels, and stiff wires for hard, dense, and blunt occlusions. To treat a narrowing in a blood vessel, a wire is passed through the stenosis in the vessel and a balloon on a catheter is passed over the wire and into the desired position. The positioning is verified by fluoroscopy and the balloon is inflated using water mixed with contrast dye to 75 to 500 times normal blood pressure (6 to 20 atmospheres), with most coronary angioplasties requiring less than 10 atmospheres. A stent may or may not also be placed.
Technique:
At the conclusion of the procedure, the balloons, wires and catheters are removed and the vessel puncture site is treated either with direct pressure or a vascular closure device.Transradial artery access (TRA) and transfemoral artery access (TFA) are two techniques for percutaneous coronary intervention. TRA is the technique of choice for management of acute coronary syndrome (ACS) as it has significantly lower incidence of bleeding and vascular complications compared with the TFA approach. TRA also has a mortality benefit for high risk ACS patients and high risk bleeding patients. TRA was also found to yield improved quality of life, as well as decreased healthcare costs and resources.
Risks and complications:
Relative to surgery, angioplasty is a lower-risk option for the treatment of the conditions for which it is used, but there are unique and potentially dangerous risks and complications associated with angioplasty: Embolization, or the launching of debris into the bloodstream Bleeding from over-inflation of a balloon catheter or the use of an inappropriately large or stiff balloon, or the presence of a calcified target vessel.
Risks and complications:
Hematoma or pseudoaneurysm formation at the access site Radiation-induced injuries (burns) from the X-rays used Contrast-induced renal injury Cerebral Hyperperfusion Syndrome leading to stroke is a serious complication of carotid artery angioplasty with stenting.Angioplasty may also provide a less durable treatment for atherosclerosis and be more prone to restenosis relative to vascular bypass or coronary artery bypass grafting. Drug-eluting balloon angioplasty has significantly less restenosis, late lumen loss and target lesion revascularization at both short term and midterm follow-up compared to uncoated balloon angioplasty for femoropopliteal arterial occlusive disease. Although angioplasty of the femoropopliteal artery with paclitaxel-coated stents and balloons significantly reduces rates of vessel restenosis and target lesion revascularization, it was also found to have increased risk of death.
Recovery:
After angioplasty, most patients are monitored overnight in the hospital, but if there are no complications, patients are sent home the following day.The catheter site is checked for bleeding and swelling and the heart rate and blood pressure are monitored to detect late rupture and hemorrhage. Post-procedure protocol also involves monitoring urinary output, cardiac symptoms, pain and other signs of systemic problems. Usually, patients receive medication that will relax them to protect the arteries against spasms. Patients are typically able to walk within two to six hours following the procedure and return to their normal routine by the following week.Angioplasty recovery consists of avoiding physical activity for several days after the procedure. Patients are advised to avoid heavy lifting and strenuous activities for a week. Patients will need to avoid physical stress or prolonged sport activities for a maximum of two weeks after a delicate balloon angioplasty.After the initial two week recovery phase, most angioplasty patients can begin to safely return to low-level exercise. A graduated exercise program is recommended whereby patients initially perform several short bouts of exercise each day, progressively increasing to one or two longer bouts of exercise. As a precaution, all structured exercise should be cleared by a cardiologist before commencing. Exercise-based rehabilitation following percutaneous coronary intervention has shown improvement in recurrent angina, total exercise time, ST-segment decline, and maximum exercise tolerance.Patients who experience swelling, bleeding or pain at the insertion site, develop fever, feel faint or weak, notice a change in temperature or color in the arm or leg that was used or have shortness of breath or chest pain should immediately seek medical advice.
Recovery:
Patients with stents are usually prescribed dual antiplatelet therapy (DAPT) which consists of a P2Y12 inhibitor, such as clopidogrel, which is taken at the same time as acetylsalicylic acid (aspirin). Dual antiplatelet therapy (DAPT) is recommended for 1 month following bare metal stent placement, for 3 months following a second generation drug-eluting stent placement, and for 6–12 months following a first generation drug-eluting stent placement. DAPT's antiplatelet properties are intended to prevent blood clots, however they also increase the risk of bleeding, so it is important to consider each patient's preferences, cardiac conditions, and bleeding risk when determining the duration of DAPT treatment. Another important consideration is that concomitant use of Clopidogrel and Proton Pump Inhibitors following coronary angiography is associated with significantly higher adverse cardiovascular complications such as major adverse cardiovascular events (MACE), stent thrombosis and myocardial infarction.
History:
Angioplasty was first described by the US interventional radiologist Charles Dotter in 1964. Dr. Dotter pioneered modern medicine with the invention of angioplasty and the catheter-delivered stent, which were first used to treat peripheral arterial disease. On January 16, 1964, Dotter percutaneously dilated a tight, localized stenosis of the subsartorial artery in an 82-year-old woman with painful leg ischemia and gangrene who refused leg amputation. After successful dilation of the stenosis with a guide wire and coaxial Teflon catheters, the circulation returned to her leg. The dilated artery stayed open until her death from pneumonia two and a half years later. Charles Dotter is commonly known as the "Father of Interventional Radiology" and was nominated for the Nobel Prize in medicine in 1978.
History:
The first percutaneous coronary angioplasty on an awake patient was performed in Zurich by the German cardiologist Andreas Gruentzig on September 16, 1977.The first percutaneous coronary angioplasties in the United States were performed on the same day (March 1, 1978) by Dr. Simon H. Stertzer at Lenox Hill Hospital in New York and Dr. Richard K. Myler at St. Mary's Hospital in San Francisco. During the previous year, also at St. Mary's Hospital in San Francisco, Drs. Myler and Gruentzig had performed dilatations in the setting of bypass surgery to test the catheter concept before Gruentzig performed the first PTCA in his catheterization lab in Zurich.
History:
The initial form of angioplasty was 'plain old balloon angioplasty' (POBA) without stenting, until the invention of bare metal stenting in the mid-1980s to prevent the abrupt closure sometimes seen with POBA.Bare metal stents were found to cause in-stent restenosis as a result of neointimal hyperplasia and stent thrombosis, which led to the invention of drug-eluting stents with anti-proliferative drugs to combat in-stent restenosis.The first coronary angioplasty with a drug delivery stent system was performed by Dr. Stertzer and Dr. Luis de la Fuente, at the Instituto Argentino de Diagnóstico y Tratamiento (English: Argentina Institute of Diagnosis and Treatment) in Buenos Aires, in 1999.
History:
Ingemar Henry Lundquist invented the over-the-wire balloon catheter that is now used in the majority of angioplasty procedures in the world. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Police uniforms in the United States**
Police uniforms in the United States:
Police uniforms in the United States vary widely due to the nation's tradition of highly decentralized law enforcement. Over time, however, a number of general conventions and styles have become representative of American police fashion. Police officers wear uniforms to deter crime by establishing a visible presence while on patrol, to make themselves easily identifiable to non-police officers or to their colleagues who require assistance, and to quickly identify each other at crime scenes for ease of coordination.
History:
Centralized, municipally-managed police departments were unknown in the United States prior to the 1830s. Early law enforcement functions were largely performed by volunteer watchmen as well as elected or appointed constables and sheriffs, who were paid by the fee system for warrants they served. The advent of professional police forces in the United States foreshadowed the introduction of standardized police uniforms. While uniforms for police had been introduced in the United Kingdom as early as 1828, adoption of standardized dress in the United States took longer, with many of the new police objecting to uniforms out of concern they would be subject to public ridicule. Nonetheless, in 1854, the New York City Police Department became the United States' first municipal police force to issue uniforms to its officers. New York City was followed, in 1858, by Boston, Chicago, and soon thereafter, other cities.
History:
The navy blue uniforms adopted by many police departments in this early period were simply surplus United States Army uniforms from the Civil War. Headwear typically took the form of stovepipe hats, a starched woolen head cover similar in appearance to a top hat but with a squatter dimension, or British-style custodian helmets. In rural areas, where preventative policing was limited or non-existent, sheriff's deputies continued to wear civilian attire, using only their badge as a mark of identification. In many states this practice continued well into the following century. The Orange County, California, sheriff's office, for instance, did not adopt a uniform until 1938.By the early 20th century, the style and form of American police uniforms had essentially settled into its modern pattern of collared shirts, neckties, slacks and military-style jackets with open collars, all worn with peaked hats. Many early uniforms had loose-fitting jackets that would conceal a police officer's equipment, such as truncheon and sidearm. Beginning in the 1930s, officers more frequently began wearing their personal gear on a Sam Browne belt worn outside the coat, for ease of access.
History:
One of the biggest evolutionary experiments in police uniform design began in 1969, when the police department in Menlo Park, California moved away from typical police uniforms, opting instead for a dress style designed to better emulate civilian fashion trends and communicate a "softer" appearance. The new uniforms consisted of green blazers, black slacks, a white shirt and black necktie. Officers wore their weapons concealed under their coats. Many other police departments soon followed the Menlo Park lead. In psychological tests, it was discovered police - after using the new uniform - displayed less authoritarian personality characteristics. In addition, civilians and suspects injured during arrests by police dropped by 50-percent and assaults on officers by suspects also plummeted by nearly a third. Despite these initially promising signs, however, it was subsequently determined that other factors, including increased police recruitment of college graduates and adoption of more responsive management techniques, had probably accounted for the statistical shifts. By the eighth year of the uniform experiment, assaults on police had more than doubled from what they were prior to the dress change and the "civilian" style uniforms were subsequently dropped.
Current designs:
Badges Despite the wide variety of uniforms used by United States police departments, virtually all incorporate the use of metallic badges as a means of primary identification. Unlike in the United Kingdom, where officers both in and out of uniform carry - but do not publicly display - paper or plastic warrant cards, US police badges are the official symbol of office and are prominently worn over the left chest of the uniform (or, in the case of plainclothes officers, displayed from a concealed badge carrier when necessary to establish authority). In Virginia, for instance, police only have the power to make arrests when "in uniform, or displaying a badge of office."Badges are typically engraved with a unique identification number matched to the officer to whom it is issued. Some departments - most notably the New York City Police Department (NYPD) - traditionally pass individual badges through several generations of police so that current officers can establish a symbolic connection with the retired and deceased officers to whom their badge had previously been issued. In the case of the NYPD, officers who misplace their badge are docked five days of vacation time and many officers wear replica badges to avoid losing their issued badge (though the practice is officially discouraged).Federal law prohibits the sale or purchase of counterfeit police badges and many states have laws regulating the wearing of metallic badges by persons other than law enforcement. Florida, for instance, prohibits unauthorized persons from wearing or displaying badges if their wear or display would be likely to deceive someone. New York, Connecticut, Massachusetts, Maine and New Jersey, meanwhile, allow private security guards to wear badges provided they are in the shape of a square and not the more traditional shield or star shape used by police.
Current designs:
Badges are usually constructed out of metal with an enamel finish in either a gold and/or silver. As a general rule, the badges issued by county sheriff's offices take the form of a five, six, or seven-pointed star, while municipal police have shield-like designs. Following the death of a police officer, other officers will typically cover their badges with a black mourning band. Mourning bands can also be seen worn on May 15, "National Peace Officers Memorial Day." Patches Most police uniforms feature shoulder sleeve insignia in the form of cloth patches embroidered with the agency's name, logo or a heraldic device. These patches are displayed either on both shoulders, or one.
Current designs:
Uniforms Individual municipal and county law enforcement agencies in the United States are typically responsible for designing their own uniforms, often with minimal state regulation. As a result, there is no universal form or pattern for American police uniforms.
Current designs:
However, in general, most large police departments provide officers with two types of uniforms for wear, tactical (also called "Class B"), and traditional (or "Class A"). Tactical uniforms, similar in material and cut to the U.S. Army's former battle dress uniform, are generally worn while on patrol, or performing physically intense duties, while traditional-style uniforms are more often used for station assignments, high-profile events, and ceremonial functions. In addition to these two basic uniform types, a variety of specialized clothing may be deployed as necessary, including jumpsuits (sometimes called "Class C") and, in the case of police pipe bands, highland dress. Many police departments restrict the use of tactical uniforms to tactical units, such as SWAT teams, or for special assignments, such as riot control, in order to present a less militarized appearance in day-to-day operations.
Current designs:
Municipal police uniforms are typically colored in blue or black, while uniforms worn by sheriff's deputies are more often green, brown, or khaki. Unlike British police, American law enforcement agencies do not usually include white-colored apparel, such as shirts, in their uniforms due to the fact white reflects in the dark and can make police officers more prominent targets for armed criminals during building searches or standoffs. There are, however, many exceptions to this general rule; the Miami Police Department wears dark blue, bike division officers will wear white, and they may also wear polo shirts on some occasions, and senior officers in the New York City Police Department, Baltimore Police Department, Philadelphia Police Department, and Washington, D.C. Metropolitan Police Department wear uniforms that also feature white shirts. However, some departments like the Hillsborough County Sheriff's Office have white shirts as part of the universal standard uniform.
Current designs:
Headgear Unlike police in some Commonwealth nations, US police forces subsequently abandoned the custodian helmet by the early 20th century they had incorporated with their uniforms. Today, municipal police forces typically wear peaked hats or, in tactical uniforms, baseball caps. County sheriff's offices often issues their deputies with campaign hats or Stetsons for cover. Some departments permit the usage of the hijab for female Muslim police officers.
Current designs:
Unique elements Several United States police forces are known for unique uniform items not commonly used by other departments. Police uniforms in Chicago and Pittsburgh feature peaked hats incorporating the Sillitoe tartan checkerboard design, similar to taxicab decor.
The Washington State Patrol and New Mexico State Police wear bow ties.
Troopers of the Texas Ranger Division and Texas Highway Patrol wear cowboy hats as part of their duty and dress uniforms and cowboy boots as part of their dress uniforms.
The Honolulu Police Department uses Kukui nut medallions on their epaulets for the ranks of Lieutenant through Major.
Many police departments replace their standard collared shirts and neckties with an ascot when detachments are organized as part of a funeral detail or color guard. Fourragère and aiguillettes may also be worn.
Awards:
Since 1977, the North American Association of Uniforms, Manufacturers and Distributors have sponsored an annual award for the "best dressed" police department in the United States. The awards are currently awarded in several classes, depending on a department's size. In 2013 the Florida Highway Patrol was recognized as the "best dressed" police force among departments with more than 2,000 personnel. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Muffin tin**
Muffin tin:
A muffin or cupcake tray is a mold in which muffins or cupcakes are baked. A single cup within a regular muffin tin is 100 millilitres (3.5 US fl oz) and most often has room for 12 muffins, although tins holding 6, 8, 11, 24, and 35 muffins do exist. A single cup within a mini muffin tin is 62.8 millilitres (2.125 US fl oz), and because these are less common, there are several standard numbers of cups per tin, including 6, 12, and 24 cups per tin. A single cup within a jumbo muffin tin is 242.13 millilitres (8.1875 US fl oz), and again because these are uncommon, there are several standard numbers of cups per tin, including 4, 6, and 12 cups per tin.Muffin tins can be made out of aluminum, stainless steel, cast iron, or silicone. In addition, aluminum and stainless steel muffin tins may be coated with Teflon or other non-stick coatings. Historically, galvanized steel has been used for muffin tins but this is no longer common. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Debian Conference**
Debian Conference:
DebConf, the Debian developers conference is the yearly conference where developers of the Debian operating system meet to discuss further development of the system.
Besides the scheduled workshops and talks, Debian developers take the opportunity to hack on the Debian system in a more informal setting. This has been institutionalised by introducing the DebCamp in the Oslo DebConf in 2003: a room is set aside and computing infrastructure provided.
Locations:
Locations of past and future DebConf events:
Miniconf:
These were one-day miniature conferences, originally held in association with the main linux.conf.au Australian Linux conference. They were targeted towards specific communities of interest and offered delegates an opportunity to network with other enthusiasts while immersing themselves in a specific topic or project.
Locations of past LCA Miniconf events:
MiniDebConf:
This is a smaller Debian event, held annually in various places in the world.
Locations of past and future MiniDebConf events:
Attendance:
According to a 2013 brochure, the conference had about 30 attendees in 2000 while in 2011 there were around 300 attendees, and about 250 are expected. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Southerly Buster**
Southerly Buster:
A southerly buster is the colloquial name of an abrupt southerly wind change in the southern regions of New South Wales and Victoria, Australia, which approaches from the southeast, mainly on a hot day, bringing in cool, usually severe weather and a dramatic temperature drop, thus ultimately replacing and relieving the prior hot conditions. Marking the boundary between hot and cool air masses, a southerly buster is sometimes represented by a roll-up cloud perpendicular to the coast, which appears from the south and coexists with the wind change, though sometimes there is little visual signal of the southerly's arrival.Southerly busters occur in the backside of a low pressure trough, followed by the speedy advance of an anticyclone from Western Australia. They are caused by the interaction of a shallow cold front with the blocking mountain range that aligns the coast, and frictional contrasts over the mainland and the ocean that disconnect the flow. Southerly busters occur about 32 times each year on the coast of southeastern Australia, with variable strength, usually in spring and summer. Although southerly busters are often associated with NSW and Victoria, they also occur on the east coast of Tasmania, New Zealand, and in Argentina and Chile.
History:
19th century In 1819, English explorer William Wentworth describes the southerly buster experienced in Sydney: ...They [hot northwesterly winds] seldom, however, continue for more than two days at a time, and are always superseded by a cold southerly gale, generally accompanied with rain. The thermometer then sinks sometimes as low as 60 °F (16 °C), and a variation of temperature of from 30 to 40 °F (17 to 22 °C) takes place in half an hour. These southerly gales usually last at this season from twelve to twenty-four hours, and then give way to the regular sea and land breezes...
History:
From the early days of settlement at Sydney Cove, sudden southerly squalls had been a problem for boats in the harbour. In 1829, a brickfielder as these squalls as were then called, laid the boat carrying Governor Ralph Darling and his family on its beam ends. Later Port Jackson boatmen would call it, the "Sútherly Búster". During the early days of the European settlement, Sydney's summer storms were accompanied by rolling red dust from the colony's brickworks.In 1869, following the connection of lighthouses and signal stations on the east coast of NSW to the electric telegraph network, a system was instituted to give shipping warning of approaching gales. It would give boats in Port Jackson warning of approaching busters. The new signals were included in Gowland's amendments to the NSW Sailing Directory The signal was a diamond shape (see Day shapes) on the southern yard-arm of the signal mast to indicate a squall approaching from the south, and numeral pendants per John Nicholson's Code of Signals flying on the masthead to show the location of the squall. The numbers were those already in use from 1842 for identifying ports, e.g. Jervis Bay was 46 and Wollongong, added in 1854, was 82. The signal masts were at South Head Signal Station and Fort Phillip. These masts had two yard-arms at right angles aligned north–south and east–west. The success of the system was seen in 1870 when a newspaper reported, "About 9 o'clock, however, intelligence of a southerly burster was telegraphed from Eden, and by half-past 12 the expected 'brickfielder' blew over the city."From 1876 signals would be displayed at two additional places - the new flagstaff at Bradley's Head and the Floating Light Vessel at the entrance to the Harbour. This direction specified that a cone signal on the flagstaff would indicate the approach of a southerly squall. This signal, "A cone, with the point downwards, shows that a gale is probable; at first from the southward.", was introduced in England 1861 by Robert Fitzroy, but this appears to be the first time it was used in Port Jackson. In 1879, George Herbert Gibson published a book called Southerly Busters, where a buster wreaks havoc near Hyde Park, Sydney.
History:
In November 1898, the Postmaster General, Varney Parkes, announced that a white flag with the letters E, J.B. or W would be flown from the flagstaff on the clock tower of the General Post Office, Sydney to signal the approach of a southerly buster. As the wind reached Eden, Jervis Bay or Wollongong the relevant flag would be hoisted. The J.B. flag was hoisted for the first time on 10 February 1899 but the wind did not travel beyond Kiama. The day had been hot and oppressive and the much anticipated cool change didn't eventuate. The original intention was to change the flags as the change moved north but this didn't always happen. It was possibly later reduced to just two flags, JB and W.
History:
20th century On the night of 16 December 1908 a new signal, also on the GPO clock tower, made its debut. For warnings at night a red light had been installed in the lantern - the glass enclosed room below the flagstaff. It continued until 1940. Capsizes still occurred, but fishing boats did heed the warnings.In 1931, with the addition of signals for southerly busters, the practice of displaying the Bureau of Meteorology weather forecast with flags was extended, with the cooperation of the Royal Australian Navy, to the signal staff on Garden Island. The weather flags had been flown since 1912 from the flagstaffs on prominent buildings but as well, at Garden Island, a diamond shape would be hoisted for warnings during the day and a red light at night when a southerly buster was expected. In 1932 Navy League sea cadets flew the same signals as at Garden Island at their station on Snapper Island (New South Wales).The Weather beacon on the top of MLC Building, North Sydney, built in 1957, could also warn of 'southerly busters'. Generally, the forecasts displayed by the lights were updated three times a day but could be changed at anytime if a sudden change was imminent. If a 'buster' was expected the red lights at the bottom of the beacon would flash at half second intervals. This signal, indicating strong winds and rain, was not specific to 'southerly busters', but as the jingle advised: Flashes short, prepare for gales, Gather the washing, furl the sails." Author Ruth Park makes a reference to the southerly buster in her novel Poor Man's Orange (1949): After an unbearably hot day, the old men on the balconies were sniffing the air and saying, 'Here she comes!' The Southerly Buster, the genie of Sydney, flapped its coarse wing over the city ... The women undid the fronts of their frocks and the little children lifted up their shirts and let it blow on their sweaty bottoms.
Formation:
On a hot day, a strong offsea gale develops from the south usually in the late afternoon and early evening, causing a rapid fall in temperature as it arrives, and sometimes a short rain and/or thunderstorm may accompany, especially if it is affiliated with a cold front coming from the south and a trough, with the strongest winds being at the leading edge of the buster. The southerly buster, which usually trails a thick layer of low stratus clouds bringing episodic drizzle, banks up against the Great Dividing Range, thus resulting in the transmission of the cool maritime air near the southeastern part of the ranges and the blocking of the airflow on the mountains' west side. The Buster progresses into a strongly stable boundary layer with warm, prefrontal landmass air that is between 100 and 200 metres deep transporting over the cooler waters.In some occasions, a cold front in the Southern Ocean may interact with the Great Dividing Range and develop into a Southerly Buster. It is worth noting that some Southerly Busters are not southern ocean fronts, as they have developed on the pre-frontal trough south of Australia or have possibly originated on the southern NSW coast in affiliation with a shallow mesoscale low development. As the Southerly Buster advances to the north, its force would decrease, though in rare occasions it can fortify north of the Hunter Valley due to upslope motion on the valley's north, which allows the flow to produce anticyclonic vorticity akin to the original obstruction and arrangement of the Southerly Buster in Victoria.However, the southerly buster does not always create precipitation, aside from light drizzle and light rain, which tend to occur a day after the southerly buster's arrival as its effects may still persist for 24 hours, in addition to creating a weather pattern similar to that of a June Gloom experienced in Southern California. It is proposed that the southerly buster is basically a coastal gravity current that is held against the mountains by the Coriolis force and in transverse geostrophic balance, and is generated when a cold front is obstructed, experiencing anticyclonic distortion near the Great Dividing Range, spreading northward as a coastal trapped orographic jet.The southerly buster is caused by its interaction with the Great Dividing Range, as the cool air becomes trapped against the ranges, oftentimes in the Gippsland area of Victoria, where the mountains create a channelling effect as the southerly winds move across the New South Wales coast. When the inland portion of the cold front is held against the mountains, the part over the sea proceeds to move along the shore, twisting the front into an 'S' shape. This activity continues on the southern New South Wales coast, while areas leading the front are still experiencing hot northwesterly winds. Other phenomena that lead to the temperature gradient between the warm air mass and the cold density current include; hot north-westerly or warm dry foehn wind that precede the squall. Furthermore, severe thunderstorms may come from the forced elevation of warm, humid air.
Formation:
Characteristics The main distinguishing feature of a southerly buster is the sudden, squally southerly change in wind direction which replace the continental northwesterly winds. This is accompanied by a marked temperature drop and sea level pressure rise. Wind gusts in excess of 40 knots (74 km/h) near ground level averages about three per year, which usually come about after very hot days and would tremendously ease within 30 to 60 minutes after the Buster's arrival, becoming rather light within a few hours. A regular southerly buster is between 20 and 60 nautical miles wide, with the heavier winds concentrated on the seaward strip, with its depth being around 1000m, restricted by the height of the mountains to west.An orographic jet, the southerly buster is clearly a mesoscale phenomenon, developing on about one day. Because busters seldom keep a staunch speed while advancing along the coast, its arrival has always been difficult to foretell, though meteorologists nowadays have the gain of satellite imagery and weather radar to foresee it, with wind warning issued by the Bureau of Meteorology. Temperature changes can be dramatic, with falls of 10 to 15 °C (18 to 27 °F) often occurring in a few minutes. In extreme conditions, a southerly buster may lower the temperatures from 40 °C (104 °F) to 19 °C (66 °F). To note, some southerly busters can be mild and not very pronounced, where they would arrive on lukewarm days and even during sultry thunderstorm events, bringing in light, though still noticeably cooler winds in the evening, with its affects still remaining in the following few days as well in some cases.
Regions:
In New South Wales, southerly busters generally reach their maximum intensity between Nowra in the South Coast and Newcastle. Southerly busters rarely pass beyond the Mid North Coast or Port Macquarie – When they do they are usually strengthened by the presence of a tropical cyclone off the north coast of the state. Moreover, they are less pronounced in that region with narrower temperature changes. Not restricted to the immediate coast, southerly busters do generally impinge upon the inland regions of NSW as well, namely the regions in the southeastern block of the state, such as, the Greater Western Sydney area, eastern portion of the Central Tablelands region going eastwards from Bathurst, the lower Hunter Valley region, Albury in the South West Slopes in a few occasions, Cooma in the Monaro (New South Wales) to the south, Canberra in the ACT, Goulburn in the Southern Tablelands and the Snowy Mountains, with the latter regions having the most dramatic cool changes due to their southerly vicinity. Sydney receives an average of about five Southerly Busters a year, mostly in late spring and early summer, with the stronger ones generally reaching the city in the late afternoon or early evening on a hot day, though at times it would arrive after several days of hot weather. It is a crucial weather feature in the Sydney area, particularly for yachtsmen.
Regions:
In Victoria, most notably in Melbourne, southerly busters occur during the afternoons where the domineering heat, brought on by north-westerly winds from central Australia, suddenly gives way to a rapid drop in temperature, followed by rain, thunderstorms and a relatively cool night. They would reach as far inland as Swan Hill in the north and Omeo to the east, but would be less pronounced and intense as one moves more inland to the north. Temperature drops in these parts of Victoria are more dramatic than those in the east coast of New South Wales, where a 10 °C (50 °F) drop can occur within half an hour (part of the easily changeable weather). Southerly busters most emerge in spring, as the landmass northbound of Melbourne starts to warm up. Meanwhile, though, the Southern Ocean, which provides cool breezes to Victoria from the west, does not warm up as swiftly as the mainland. As such, the temperature difference between hot air from the north and cold from the ocean would be very great, thus providing good conditions for the formation of thunderstorms.
Regions:
In South Australia, Mount Gambier would the most affected by southerly busters in the state due to its southeasterly location. The buster may also reach Adelaide and Whyalla in some occasions.
In Western Australia, the Buster occasionally reaches Eucla due to its location near the Southern Ocean in the warm months, where it can experience dramatic temperature drops.
In Wellington, New Zealand these storms are normally short and frequently have winds gusting between 120 km/h and 160 km/h though higher speeds are known.
In South America, these southerly fronts frequently encroach on the southern coast of Chile and Argentina and would then advance northwards on both sides of the Andes.
Records The strongest recorded Southerly Buster fell around 4:30 pm on 12 March 2010 in Wellington, New Zealand, with a maximum gust recorded of 72 knots (133 km/h; 83 mph).
Regions:
Henry Ambrose Hunt in "An Essay on Southerly Bursters" included the Dandenong Gale in his list of bursters and wrote "the wind attained, locally, to the abnormal velocity of one hundred and fifty three miles per hour (246 km/h) in a gust, the rate of one hundred and twelve miles for ten minutes and fifty seven miles per hour for nine hours." The sail-steamer Dandenong foundered off Jervis Bay during the gale with the loss of forty lives. - hence the name of the gale.
Regions:
This velocity is no longer officially recognised but the 'Kurnell Tornado' in December 2015 shows gusts in excess of 200 km/h (120 mph) are possible in the area, and the day before the tornado the possibility of a southerly buster on the following day had been proposed.
Incidents:
On 21 November 2016, at around 6pm, a powerful southerly change occurred in Melbourne, which resulted in the death of 10 people, who were asthmatic and succumbed to respiratory failure. Thousands of others across the city experienced allergic reactions and asthma-like symptoms triggered by the storm. This was due to a stark southerly wind (60 km/hour) that distributed ryegrass pollen into the moist air, rupturing them into very fine specks, small enough particles to enter people's lungs, as they were sucked up into the warm updraft of air forming the storm cells, before they returned to earth in the storm's cool down-draft, spreading across the land in the storm's efflux area. Hospitals and medical centres in the city had to arduously manage 8,500 emergency calls in the space of just five hours, and the hospitalisation of 1400 people. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Biofactories**
Biofactories:
The name biofactory comes from the improvements which different installments of traditional health services (wastewater treatment and water treatment plants) have been through; this has been done by reassessing them into a Circular Economy (CE). The concept was first used by Chilean company Aguas Andinas’ CEO, Narcís Berberana. Yves Lesty, Carlos Castro and Lisette Gajardo were the engineers that took part in its origin.
Biofactories:
Overall, biofactories have gone from a lineal processing approach, where resources are extracted and then processed, thus generating waste; to factories that supply with new valuable resources, such as electricity, natural gas, agricultural fertilizer or clean water, which are obtained from what it used to be considered as waste.
The first biofactory complex was created in 2017 by Aguas Andinas, when this new strategy was applied to its different wastewater treatment plants, La Farfana and Mapocho-Trebal; and La Florida water treatment plant, grouping them all together under the name of Biofactoría Gran Santiago.
During its first year, Biofactoría Gran Santiago generated a total of 51.792.240 kWh of electrical energy, 160.337 BTU of biogas, 111.842 tonnes of fertilizer destined to agricultural crops, and 603 millions cubic metre of clean water that was returned to its natural sources.
Principles:
According to the company that coined this new concept, biofactories operate under 6 main principles: Water run: Using the energy of water as a resource, they create new opportunities for its use after being cleaned; such as irrigation and underwater aquifer refill.
Run with its own energy: Using a renewable source such as water, the main objective of a biofactory is to take over the total amount of its operations’ electrical requirements and feed its own system, both as energy and as biomethane.
Waste transformation: They understand waste as a source of new resources, such as the mud that comes from the process of wastewater cleaning, which is used as agricultural fertilizer.
Air protection: Within its facilities, they have deodorization systems that neutralize the odour, they reduce gas emission, and they take over the reforestation of its surroundings; minimizing its environmental impact.
Biodiversity protection: They promote initiatives that generate an optimal ecosystem for wildlife. A few of them are the cleansing of rivers contaminated with wastewater, the preservation of its environment, and the creation of artificial lakes that aim to protect the same environment.
Social involvement: Community engagement: Funding of social initiatives.
Collaboration with other companies: Boosting of projects that deal with the processing of water that comes from other areas.
Research: Enhancing both biotechnology and process engineering investigations.
Acknowledgements:
In September 2018, Aguas Andinas’ biofactory project was awarded with the Momentum for Change recognition, given by the United Nations (UN) during its Climate Change conference; acknowledging its contribution to this new strategy that seeks to prevent climate change.
Furthermore, in October 2018, Chilean foundation Recyclápolis gave its Recyclápolis Environmental National Award to Biofactoría Gran Santiago's biomethanation plant, implemented by Chilean companies Metrogás and Aguas Andinas. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Germ-Soma Differentiation**
Germ-Soma Differentiation:
Germ-Soma Differentiation is the process by which organisms develop distinct germline and somatic cells. The development of cell differentiation has been one of the critical aspects of the evolution of multicellularity and sexual reproduction in organisms. Multicellularity has evolved upwards of 25 times, and due to this there is great possibility that multiple factors have shaped the differentiation of cells. There are three general types of cells: germ cells, somatic cells, and stem cells. Germ cells lead to the production of gametes, while somatic cells perform all other functions within the body. Within the broad category of somatic cells, there is further specialization as cells become specified to certain tissues and functions. In addition, stem cell are undifferentiated cells which can develop into a specialized cell and are the earliest type of cell in a cell lineage. Due to the differentiation in function, somatic cells are found ony in multicellular organisms, as in unicellular ones the purposes of somatic and germ cells are consolidated in one cell.
Germ-Soma Differentiation:
All organisms with germ-soma differentiation are eukaryotic, and represent an added level of specialization to multicellular organisms. Pure germ-soma differentiation has developed in a select number of eukaryotes (called Weismannists), included in this category are vertebrates and arthropods- however land plants, green algae, red algae, brown algae, and fungi have partial differentiation. While a significant portion of organisms with germ-soma differentiation are asexual, this distinction has been imperative in the development of sexual reproduction; the specialization of certain cells into germ cells is fundamental for meiosis and recombination.
Weismann barrier:
The strict division between somatic and germ cells is called the Weismann barrier, in which genetic information passed onto offspring is found only in germ cells. This occurs only in select organisms, however some without a Weismann barrier do present germ-soma differentiation. These organisms include land plants, many algaes, invertebrates, and fungi whose germ cells are derived from prior somatic cells as opposed to stem cells. The Weismann barrier is essential to the concept of an immortal germline, which passes down genetic information through designated germ cells.
Weismann barrier:
Organisms with germ-soma differentiation but no Weismann barrier often reproduce through somatic embryogenesis.
Benefits and Detriments of Differentiation:
There is no single widely accepted theory on the origins of somatic-germline differentiation, however of those that do exist many are based on the evolutionary advantage of differentiated multicellularity which has allowed it to survive. These theories include the development of colonial organization structures in which the division of labor between cells allowed for improvements in fitness.
Benefits and Detriments of Differentiation:
The division of labor within multicellular organisms can offer significant advantages over unicellular counterparts. Division can allow organisms to become larger, or interact with the environments (and thus fill different niches) that increase fitness. In addition to internal benefits, there is evidence that these also improve defenses against predation. On the other hand, multicellularity comes with increased energy use devoted to maintaining homeostasis instead of to reproduction.
Benefits and Detriments of Differentiation:
Dirty Work Hypothesis One major theory as to the proliferation of organisms with cell differentiation is the dirty work hypothesis. This hypothesis posits that when an organism has differentiated cells, somatic cells are able to devote energy solely to maintaining homeostasis instead of reproduction while germ cells do the opposite. One reason proposed for the relative success of the "dirty work" system of organization is that it helps manage the detrimental effects of metabolic activity, and allow for more efficient energy distribution throughout an organism. The other major reason proposed is that it prevents metabolic activity within the cell from damaging genetic material. Said activity in mitochondria and chloroplasts creates mutagenic byproducts, so in organisms with differentiation where germ cells do not engage in metabolic activity the germline is not impacted.
Uncertainty:
Due to the nature of research around the origin of life and multicellularity, it has been difficult to obtain a case study that is optimal for observing somatic-germline differentiation. One case that has been extensively studied is that of organisms in the Volvocacaeae family. Within volvocavea, there is much diversity in organizational structure, with some organisms being unicellular, colonial, or (arguably) multicellular. Within volvocine algae three genes have been identified as crucial to the development of soma cells which regulate coding for asymmetric division of cells, preventing reproductive development of soma cells, and preventing the development of somatic characteristics in germ cells (such as those meant for mobility or metabolic activity). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**TRIANGLE disease**
TRIANGLE disease:
TRIANGLE disease is a rare genetic disorder of the immune system. TRIANGLE stands for “TPPII-related immunodeficiency, autoimmunity, and neurodevelopmental delay with impaired glycolysis and lysosomal expansion” where TPP2 is the causative gene. This disease manifests as recurrent infection, autoimmunity, and neurodevelopmental delay. TRIANGLE disease was first described in a collaborative study by Dr. Helen C. Su from the National Institute of Allergy and Infectious Diseases, National Institutes of Health, and Dr. Sophie Hambleton from the University of Newcastle and their collaborators in 2014. The disease was also described by the group of Ehl et al.
Genetics:
TRIANGLE disease is caused by loss-of-function mutations in the gene TPP2, which stands for tripeptidyl peptidase II. TPP2 maps to human chromosome 13q32-q33, has 32 exons, and encodes for a 1249 amino acid protein. The genetic model for this disease is loss-of-function. This means that for people with TRIANGLE disease, the gene TPP2 is unable to produce protein or produce functional protein.Functionally, TPPII has a key role in recycling amino acids, which are protein building blocks, a fundamental cellular process. Although the body can use alternative amino acid recycling pathways to compensate for loss of TPPII, the up-regulation of alternative pathways can cause new cellular abnormalities in itself with subsequent effects on glycolysis, adaptive immunity, and innate immunity. Consequently, individuals without functioning TPPII have severe disease.
Genetics:
Inheritance TRIANGLE disease is inherited in an autosomal recessive manner. In autosomal recessive inheritance, two copies of an abnormal gene must be present in order for the disease to develop. Typically, this means both parents of an affected child silently carry one abnormal gene. This also means this also explains why reported cases of TRIANGLE disease have involved consanguinity or geographically isolated communities.Parents of a child with TRIANGLE disease have a 25% chance of having another affected child with each pregnancy. This risk is independent of prior children's status. For example, if the first two children in a family are affected, the next child has the same 25% risk of inheriting the mutation. All affected individuals have two abnormal copies of TPP2. Children who inherit only one abnormal copy of TPP2 will not develop TRIANGLE disease although they may have affected children, particularly if they marry within the family.
Diagnosis:
Clinical manifestations Clinically, TRIANGLE disease is characterized combined immunodeficiency, severe autoimmunity, and developmental delay. Patients typically present in early childhood with recurrent bacterial and viral infections of the middle ear and respiratory tract. Additionally, patients develop severe, difficult to treat autoimmunity. This autoimmunity includes auto-antibody mediated destruction of red blood cells, neutrophils, and platelets; central nervous system lupus erythematous with stroke; and hepatitis. Patients also have mild to moderate developmental delay.
Diagnosis:
Laboratory manifestations The clinical symptoms are caused by abnormalities of the immune system and disruption of basic cellular functions. Patients show markedly decreased circulating T cells, B cells and natural killer (NK) cells, with severely reduced naive T cells and hypergammaglobulinemia.
Treatment:
Once a diagnosis is made, the treatment is based on an individual's clinical condition and may include standard management for autoimmunity and immunodeficiency. Hematopoietic stem cell transplantation has cured the immune abnormalities in one TRIANGLE patient, although the neurodevelopmental delay would likely remain. Investigators at the National Institute of Allergy and Infectious Diseases at the US National Institutes of Health currently have clinical protocols to study new approaches to the diagnosis and treatment of this disorder. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Rhoifolin**
Rhoifolin:
Rhoifolin is a chemical compound. It is first isolated from plant Rhus succedanea. The term "Rhoi" derived from generic name of plant Rhus. It is a flavone, a type of flavonoid isolated from Boehmeria nivea, China grass or ramie (leaf), from Citrus limon, Canton lemon (leaf), from Citrus x aurantium, the bigarade or bitter orange (plant), from Citrus x paradisi, the grapefruit (leaf), from Ononis campestris, the cammock (shoot) and from Sabal serratula, the serenoa or sabal fruit (plant). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Hard landscape materials**
Hard landscape materials:
The term hard landscape is used by practitioners of landscape architecture and garden design to describe the construction materials which are used to improve a landscape by design. The corresponding term soft landscape materials is used to describe vegetative materials such as plants, grasses, shrubs, trees, etc. to improve landscape or outdoor space.
Types:
A wide range of hard landscape materials can be used, such as brick, gravel, rock or stone, concrete, timber, bitumen, glass, and metals.
Common gravel types include pea gravel and crushed granite gravel.'Hard landscape' can also describe outdoor furniture and other landscape products.
Advantages:
Hard-landscape materials like gravel may avoid the need for mowing, watering, fertilizing, and pesticides. Gravel lawns may also be useful in regions where grass doesn't grow well due to insufficient sunlight. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**G&T-Seq**
G&T-Seq:
G&T-seq (short for single cell genome and transcriptome sequencing) is a novel form of single cell sequencing technique allowing one to simultaneously obtain both transcriptomic and genomic data from single cells, allowing for direct comparison of gene expression data to its corresponding genomic data in the same cell...
Background:
The advent of single-cell sequencing has provided researchers with the tools to resolve genotypically and phenotypically distinct cells within a mixed population. In cases where such heterogeneity is relevant, such as in tumours, this technique enables the study of clonal relationships and tumour evolution. As well, rare cell types and samples otherwise containing low cell numbers, such as in the case of circulating tumour cells, can also be studied in greater detail. However, previous methods of library preparation typically involve the capture of either mRNA or genomic DNA (gDNA), but not both. By simultaneously capturing and sequencing both DNA and RNA through a method called G&T sequencing, researchers are able to obtain sequence information for both genome and transcriptome analysis from single cell libraries, thereby allowing integrated studies involving both networks. As a proof of concept, the authors of G&T-seq demonstrated its ability to acquire both the messenger RNA (mRNA) and genomic DNA (gDNA) by using paramagnetic beads with biotinylated oligo-deoxy-Thymine(dT) primer to separate the polyadenylated (Poly-A) RNA from its gDNA prior to amplification and library preparation. Validation experiments on G&T-seq performed using cell lines with previous sequencing data available show that sequencing coverage, gene expression profile, and DNA copy number profiles were reliably reproduced by G&T sequencing, and that this method was able to call a majority (87%) of all previously annotated single nucleotide variants (SNVs) in these cell lines. The authors have argued on this basis that the process of physically separating mRNA from gDNA did not negatively affect the yield or quality of sequencing data.
Methods:
Similar to conventional single-cell sequencing, G&T-seq involves the harvesting and lysis of desired cells. However, both gDNA and polyA-mRNA are captured and physically separated prior to amplification and library construction for analysis using sequencing platforms.
Methods:
Separation of poly-adenylated RNA from genomic DNA G&T sequencing separates the mRNA from the gDNA using an unbiased global amplification procedure described previously. First, mRNA is isolated on specialized oligo-dT (5’-biotin-triethyleneglycol-AAGCAGTGGTATCAACGCAGAGTAC(T)30VN-3’) conjugated to streptavidin-coupled paramagnetic beads. The oligo-dT binds to the poly-A tails of processed mRNA, fishing them out from the pool of genomic material. Next, the paramagnetic beads are spatially isolated by magnetization. The genomic material remaining in the supernatant is extracted and physically separated from the mRNA.
Methods:
Amplification and Sequencing The authors that developed G&T-seq utilized and validated two methods for whole-genome amplification: Multiple displacement amplification and PicoPlex. Other methods, such as MALBAC, may be applicable but have yet to be validated.
Methods:
Multiple Displacement Amplification MDA amplification technique can be used to generate long, high quality reads that produce sequencing data of comparable quality to bulk sequencing using PCR amplification. This method involves the use of hexamer primers that bind randomly to the template, followed by DNA elongation using phi29 DNA polymerase. Upon reaching the 5’ end of a downstream primer, the polymerase displaces that elongating strand to continue synthesis. The displaced strand becomes open for pairing with more primers, allowing for amplification of the displaced strand. The process continues and produces a branched DNA library that can be cut and sequenced. The authors of the G&T technique found that, though MDA used in G&T-seq yielded genomic coverage of similar breadth as MDA performed in conventional single cell sequencing, the distribution of read coverage was less even across the genome.
Methods:
PicoPlex Though MDA produces higher quality reads suitable for SNP analysis, DNA copy number profiles generated by such a technique are not highly accurate and reproducible due to its non-uniform amplification. An alternate technique called PicoPlex, developed by Rubicon Genomics, has been shown to produce better results. Here, elongation of random primers ligated to an adapter creates a complementary strand with an adapter that, when denatured and randomly reprimed, produces a double stranded fragment with complementary adapters. Denaturation into single strands allows for the formation of hairpin loops due to the complementary nature of their adapters, creating a hairpin loop library that cannot be used for subsequent amplification, thereby preventing exponential amplification of initial bias.
Methods:
cDNA amplification Messenger RNA bound to oligo-dT is reverse transcribed into cDNA using the oligo-dT primers with the addition of Template-Switching Oligo (TSO, 5"-AAGCAGTGGTATCAACGCAGAGTACrGrG+G-3’) and Superscript II reverse transcriptase. Superscript II reverse transcriptase has additional terminal transferase activity which adds a variable number of cytosine residues to the end the 3’ terminal cDNA molecule. The overhang of 3’ cytosine residues bind to the TSO, creating an extended template. The Superscript II reverse transcriptase switches templates and continues transcribing to complete the 3’ end of the cDNA. This results in a full length cDNA containing the 5’ oligo-dT primer, cDNA transcribed from mRNA, and the 3’ universal priming site for second-strand synthesis.
Methods:
The cDNA undergoes amplification using the universal primer (5’- AAGCAGTGGTATCAACGCAGAGT-3’) for 18 cycles of PCR before it undergoes library preparation using the Nextera XT Kit from Illumina and sequencing by the Illumina HiSeq platform.
Methods:
Alternatives Techniques A similar method to G&T-seq, developed months earlier, is DR-seq (DNA and RNA sequencing). The primary difference between the two techniques is the amplification step, where DNA and polyA-RNA amplification occurs without their prior separation. DR-seq uses random priming, where primers containing a common 27-nucleotide sequence along with a variable 8-nucleotide (ad2 primers) bind to different locations on the cDNA. Despite there being multiple (50-250) primer binding sites on most cDNA, each original (i.e. not the product of amplification/in vitro transcription) cDNA molecule is usually primed only once during the initial amplification step, thus creating a single amplicon of a unique length, containing the ad2 primer on the 5' end. The 3' end contains the ad1 primer, which is the original poly-dT primer used for reverse transcription. This unique amplicon is termed the length-based identifier. Importantly, the length-based identifier is created, but not amplified by this quasilinear PCR step. The number of unique length-based identifiers for each gene can then be used to infer the number of original cDNA (and thus mRNA) molecules present for the gene, providing a method of estimating gene expression that avoids the effect of amplification bias. To further amplify the cDNA for RNA-seq, the cDNA amplicons generated by the original PCR step undergoes in vitro transcription using the T7 promoter incorporated in the ad1 primer to ensure RNA transcripts come from cDNA, not gDNA.
Methods:
Advantages of the DR-seq technique include the reduction of the possibility for contamination and RNA loss, since the extra step of DNA/RNA separation is skipped. As well, amplification bias is reduced due to the use of the aforementioned length-based identifiers. However, since DNA and polyA-RNA is not separated prior to amplification and subsequent sequencing, the exonic regions must be computationally masked, leaving only reads that originate from gDNA, in order to determine copy number. This creates issues for accurately determining copy number counts from gDNA. The authors note, though, that copy number count over large genomic regions is apparently not impacted by masking as a result because coding regions compose a relatively small portion of the genome.
Applications:
Dual genome and transcriptome sequencing allows researchers to establish high resolution correlations of genomic aberrations with alterations to levels of transcription. For example, the authors of this technique were able to detect single cells with chromosomal aneuploidies, and establish that these aneuploidies corresponded with increased or decreased overall chromosomal gene expression when there was a respective chromosomal gain (e.g. Trisomy) or loss. Subchromosomal changes could also be correlated with changes in expression of genes at affected loci. As well, the authors were able to find a fusion transcript and locate the chromosomal breakpoint in the same cell resulting in the fusion.G&T-seq also provides a strategy for establishing causative links between genotype and phenotype associations in single cells (e.g. Non-coding SNVs). While bulk sequencing of genome and transcriptome may allow one to associate a collection of genotypic features with mean expression patterns in a population of cells, it overlooks subtle or temporal differences between individual cells that may arise due to cell ecology. This presents an obstacle for researchers trying to pinpoint the genomic causes underlying transcript alterations, especially when compounded with tumour samples where heterogeneity is widespread and background genetic variation could confound relevant mutations. Conventional single cell sequencing, on the other hand, prevents one from making direct associations between mutations and changes in the transcriptome because either the DNA or the RNA is lost in the process. Traditionally, researchers would have to use other methods, such as classification based on cell markers. However, such methods of discrimination rely on the availability of specific antibodies, and provide relatively coarse discrimination compared to sequencing since expression of cell surface markers constitute only a fraction of its overall phenotypeFinally, separation of DNA from RNA paves the way for dual sequencing of the epigenome and transcriptome, two components of the cell that are intricately linked to each other. However, this would require validation with conventional single cell bisulphite sequencing to ensure separation of DNA and RNA doesn’t affect the DNA methylation status.
Considerations:
GC bias The MDA amplification has an inherent bias against repeat sequences which were underrepresented in MDA products. In the context of G&T sequencing, this results in a reduced read count as the % of GC content increases for a particular region.
Distribution of read coverage Comparing the amplification of single cell residual genomic DNA after mRNA isolation by MDA to amplification of single cell genomic DNA without mRNA isolation by MDA, showed a less evenly distributed coverage across the genome after mRNA isolation. Although there was a reduction in coverage distribution, it was not by a large proportion.
Considerations:
Exclusion of alternate RNA Isolation of mRNA by the G&T-seq technique described is only capable of capturing mRNAs which have a sufficient length poly-A tail which can be captured by the oligo-dT bait. This is not a complete representation of the mRNA present in the cell. Some mRNAs have crucial roles in phenotypic expression but do not present the standard polyA tail length due to alternative polyadenylation. Therefore, G&Ts comparison of genotype–phenotype correlation does not necessarily represent the best causal link between the two.
Considerations:
Protein expression correlation The mRNA isolation is not the only hurdle in establishing genotype–phenotype relation. It is not sufficient to use mRNA as a surrogate to total protein expression, because other RNA species exist which also play important roles in phenotypic expression. Another auxiliary technique which can bolster the claims made by G&T sequencing is a total proteome analysis by mass spectrometry, giving a better presentation of the relation between genomic changes and phenotypic presentation | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Multi-frequency receiver**
Multi-frequency receiver:
Multi-Frequency signalling, (MF), is similar to the European version, CCITT Signaling System 5, (SS5). The original format was five tones used in pairs. This later evolved to six tones. Because its six tones are used only in pairs, this signaling format is sometimes referred to as "two-out-of-five code" or "two of six." Multi-Frequency receivers have been present in US telephony at least since the late 1940s. In 1940s technology, receivers in 4XB and similar equipment used vacuum tubes. Later ones used RC filters and transistors. Digital filters became commonplace in electronic switching systems of the 1980s. For example, in 5ESS switch such jobs are done by DSPs in the Global Digital Services Unit (GDSU). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Fat binary**
Fat binary:
A fat binary (or multiarchitecture binary) is a computer executable program or library which has been expanded (or "fattened") with code native to multiple instruction sets which can consequently be run on multiple processor types. This results in a file larger than a normal one-architecture binary file, thus the name.
Fat binary:
The usual method of implementation is to include a version of the machine code for each instruction set, preceded by a single entry point with code compatible with all operating systems, which executes a jump to the appropriate section. Alternative implementations store different executables in different forks, each with its own entry point that is directly used by the operating system.
Fat binary:
The use of fat binaries is not common in operating system software; there are several alternatives to solve the same problem, such as the use of an installer program to choose an architecture-specific binary at install time (such as with Android multiple APKs), selecting an architecture-specific binary at runtime (such as with Plan 9's union directories and GNUstep's fat bundles), distributing software in source code form and compiling it in-place, or the use of a virtual machine (such as with Java) and just-in-time compilation.
Apollo:
Apollo's compound executables In 1988, Apollo Computer's Domain/OS SR10.1 introduced a new file type, "cmpexe" (compound executable), that bundled binaries for Motorola 680x0 and Apollo PRISM executables.
Apple:
Apple's fat binary A fat-binary scheme smoothed the Apple Macintosh's transition, beginning in 1994, from 68k microprocessors to PowerPC microprocessors. Many applications for the old platform ran transparently on the new platform under an evolving emulation scheme, but emulated code generally runs slower than native code. Applications released as "fat binaries" took up more storage space, but they ran at full speed on either platform. This was achieved by packaging both a 68000-compiled version and a PowerPC-compiled version of the same program into their executable files. The older 68K code (CFM-68K or classic 68K) continued to be stored in the resource fork, while the newer PowerPC code was contained in the data fork, in PEF format.Fat binaries were larger than programs supporting only the PowerPC or 68k, which led to the creation of a number of utilities that would strip out the unneeded version. In the era of small hard drives, when 80 MB hard drives were a common size, these utilities were sometimes useful, as program code was generally a large percentage of overall drive usage, and stripping the unneeded members of a fat binary would free up a significant amount of space on a hard drive.
Apple:
NeXT's/Apple's multi-architecture binaries NeXTSTEP Multi-Architecture Binaries Fat binaries were a feature of NeXT's NeXTSTEP/OPENSTEP operating system, starting with NeXTSTEP 3.1. In NeXTSTEP, they were called "Multi-Architecture Binaries". Multi-Architecture Binaries were originally intended to allow software to be compiled to run both on NeXT's Motorola 68k-based hardware and on Intel IA-32-based PCs running NeXTSTEP, with a single binary file for both platforms. It was later used to allow OPENSTEP applications to run on PCs and the various RISC platforms OPENSTEP supported. Multi-Architecture Binary files are in a special archive format, in which a single file stores one or more Mach-O subfiles for each architecture supported by the Multi-Architecture Binary. Every Multi-Architecture Binary starts with a structure (struct fat_header) containing two unsigned integers. The first integer ("magic") is used as a magic number to identify this file as a Fat Binary. The second integer (nfat_arch) defines how many Mach-O Files the archive contains (how many instances of the same program for different architectures). After this header, there are nfat_arch number of fat_arch structures (struct fat_arch). This structure defines the offset (from the start of the file) at which to find the file, the alignment, the size and the CPU type and subtype which the Mach-O binary (within the archive) is targeted at.
Apple:
The version of the GNU Compiler Collection shipped with the Developer Tools was able to cross-compile source code for the different architectures on which NeXTStep was able to run. For example, it was possible to choose the target architectures with multiple '-arch' options (with the architecture as argument). This was a convenient way to distribute a program for NeXTStep running on different architectures.
Apple:
It was also possible to create libraries (e.g. using NeXTStep's libtool) with different targeted object files.
Apple:
Mach-O and Mac OS X Apple Computer acquired NeXT in 1996 and continued to work with the OPENSTEP code. Mach-O became the native object file format in Apple's free Darwin operating system (2000) and Apple's Mac OS X (2001), and NeXT's Multi-Architecture Binaries continued to be supported by the operating system. Under Mac OS X, Multi-Architecture Binaries can be used to support multiple variants of an architecture, for instance to have different versions of 32-bit code optimized for the PowerPC G3, PowerPC G4, and PowerPC 970 generations of processors. It can also be used to support multiple architectures, such as 32-bit and 64-bit PowerPC, or PowerPC and x86, or x86-64 and ARM64.
Apple:
Apple's Universal binary In 2005, Apple announced another transition, from PowerPC processors to Intel x86 processors. Apple promoted the distribution of new applications that support both PowerPC and x86 natively by using executable files in Multi-Architecture Binary format. Apple calls such programs "Universal applications" and calls the file format "Universal binary" as perhaps a way to distinguish this new transition from the previous transition, or other uses of Multi-Architecture Binary format.
Apple:
Universal binary format was not necessary for forward migration of pre-existing native PowerPC applications; from 2006 to 2011, Apple supplied Rosetta, a PowerPC (PPC)-to-x86 dynamic binary translator, to play this role. However, Rosetta had a fairly steep performance overhead, so developers were encouraged to offer both PPC and Intel binaries, using Universal binaries. The obvious cost of Universal binary is that every installed executable file is larger, but in the years since the release of the PPC, hard-drive space has greatly outstripped executable size; while a Universal binary might be double the size of a single-platform version of the same application, free-space resources generally dwarf the code size, which becomes a minor issue. In fact, often a Universal-binary application will be smaller than two single-architecture applications because program resources can be shared rather than duplicated. If not all of the architectures are required, the lipo and ditto command-line applications can be used to remove versions from the Multi-Architecture Binary image, thereby creating what is sometimes called a thin binary.
Apple:
In addition, Multi-Architecture Binary executables can contain code for both 32-bit and 64-bit versions of PowerPC and x86, allowing applications to be shipped in a form that supports 32-bit processors but that makes use of the larger address space and wider data paths when run on 64-bit processors.
Apple:
In versions of the Xcode development environment from 2.1 through 3.2 (running on Mac OS X 10.4 through Mac OS X 10.6), Apple included utilities which allowed applications to be targeted for both Intel and PowerPC architecture; universal binaries could eventually contain up to four versions of the executable code (32-bit PowerPC, 32-bit x86, 64-bit PowerPC, and 64-bit x86). However, PowerPC support was removed from Xcode 4.0 and is therefore not available to developers running Mac OS X 10.7 or greater.
Apple:
In 2020, Apple announced another transition, this time from Intel x86 processors to Apple silicon (ARM64 architecture). To smooth the transition Apple added support for the Universal 2 binary format; Universal 2 binary files are Multi-Architecture Binary files containing both x86-64 and ARM64 executable code, allowing the binary to run natively on both 64-bit Intel and 64-bit Apple silicon. Additionally, Apple introduced Rosetta 2 dynamic binary translation for x86 to Arm64 instruction set to allow users to run applications that do not have Universal binary variants.
Apple:
Apple Fat EFI binary In 2006, Apple switched from PowerPC to Intel CPUs, and replaced Open Firmware with EFI. However, by 2008, some of their Macs used 32-bit EFI and some used 64-bit EFI. For this reason, Apple extended the EFI specification with "fat" binaries that contained both 32-bit and 64-bit EFI binaries.
CP/M and DOS:
Combined COM-style binaries for CP/M-80 and DOS CP/M-80, MP/M-80, Concurrent CP/M, CP/M Plus, Personal CP/M-80, SCP and MSX-DOS executables for the Intel 8080 (and Z80) processor families use the same .COM file extension as DOS-compatible operating systems for Intel 8086 binaries. In both cases programs are loaded at offset +100h and executed by jumping to the first byte in the file. As the opcodes of the two processor families are not compatible, attempting to start a program under the wrong operating system leads to incorrect and unpredictable behaviour.
CP/M and DOS:
In order to avoid this, some methods have been devised to build fat binaries which contain both a CP/M-80 and a DOS program, preceded by initial code which is interpreted correctly on both platforms. The methods either combine two fully functional programs each built for their corresponding environment, or add stubs which cause the program to exit gracefully if started on the wrong processor. For this to work, the first few instructions (sometimes also called gadget headers) in the .COM file have to be valid code for both 8086 and 8080 processors, which would cause the processors to branch into different locations within the code.
CP/M and DOS:
For example, the utilities in Simeon Cran's emulator MyZ80 start with the opcode sequence EBh, 52h, EBh. An 8086 sees this as a jump and reads its next instruction from offset +154h whereas an 8080 or compatible processor goes straight through and reads its next instruction from +103h.
A similar sequence used for this purpose is EBh, 03h, C3h.
CP/M and DOS:
John C. Elliott's FATBIN is a utility to combine a CP/M-80 and a DOS .COM file into one executable. His derivative of the original PMsfx modifies archives created by Yoshihiko Mino's PMarc to be self-extractable under both, CP/M-80 and DOS, starting with EBh, 18h, 2Dh, 70h, 6Dh, 73h, 2Dh to also include the "-pms-" signature for self-extracting PMA archives, thereby also representing a form of executable ASCII code.
CP/M and DOS:
Another method to keep a DOS-compatible operating system from erroneously executing .COM programs for CP/M-80 and MSX-DOS machines is to start the 8080 code with C3h, 03h, 01h, which is decoded as a "RET" instruction by x86 processors, thereby gracefully exiting the program, while it will be decoded as "JP 103h" instruction by 8080 processors and simply jump to the next instruction in the program. Similar, the CP/M assembler Z80ASM+ by SLR Systems would display an error message when erroneously run on DOS.Some CP/M-80 3.0 .COM files may have one or more RSX overlays attached to them by GENCOM. If so, they start with an extra 256-byte header (one page). In order to indicate this, the first byte in the header is set to magic byte C9h, which works both as a signature identifying this type of COM file to the CP/M 3.0 executable loader, as well as a "RET" instruction for 8080-compatible processors which leads to a graceful exit if the file is executed under older versions of CP/M-80.C9h is never appropriate as the first byte of a program for any x86 processor (it has different meanings for different generations, but is never a meaningful first byte); the executable loader in some versions of DOS rejects COM files that start with C9h, avoiding incorrect operation.
CP/M and DOS:
Similar overlapping code sequences have also been devised for combined Z80/6502, 8086/68000 or x86/MIPS/ARM binaries.
CP/M and DOS:
Combined binaries for CP/M-86 and DOS CP/M-86 and DOS do not share a common file extension for executables. Thus, it is not normally possible to confuse executables. However, early versions of DOS had so much in common with CP/M in terms of its architecture that some early DOS programs were developed to share binaries containing executable code. One program known to do this was WordStar 3.2x, which used identical overlay files in their ports for CP/M-86 and MS-DOS, and used dynamically fixed-up code to adapt to the differing calling conventions of these operating systems at runtime.Digital Research's GSX for CP/M-86 and DOS also shares binary identical 16-bit drivers.
CP/M and DOS:
Combined COM and SYS files DOS device drivers (typically with file extension .SYS) start with a file header whose first four bytes are FFFFFFFFh by convention, although this is not a requirement. This is fixed up dynamically by the operating system when the driver loads (typically in the DOS BIOS when it executes DEVICE statements in CONFIG.SYS). Since DOS does not reject files with a .COM extension to be loaded per DEVICE and does not test for FFFFFFFFh, it is possible to combine a COM program and a device driver into the same file by placing a jump instruction to the entry point of the embedded COM program within the first four bytes of the file (three bytes are usually sufficient). If the embedded program and the device driver sections share a common portion of code, or data, it is necessary for the code to deal with being loaded at offset +0100h as a .COM style program, and at +0000h as a device driver. For shared code loaded at the "wrong" offset but not designed to be position-independent, this requires an internal address fix-up similar to what would otherwise already have been carried out by a relocating loader, except for that in this case it has to be done by the loaded program itself; this is similar to the situation with self-relocating drivers but with the program already loaded at the target location by the operating system's loader.
CP/M and DOS:
Crash-protected system files Under DOS, some files, by convention, have file extensions which do not reflect their actual file type. For example, COUNTRY.SYS is not a DOS device driver, but a binary NLS database file for use with the CONFIG.SYS COUNTRY directive and the NLSFUNC driver. The PC DOS and DR-DOS system files IBMBIO.COM and IBMDOS.COM are special binary images loaded by bootstrap loaders, not COM-style programs. Trying to load COUNTRY.SYS with a DEVICE statement or executing IBMBIO.COM or IBMDOS.COM at the command prompt will cause unpredictable results.It is sometimes possible to avoid this by utilizing techniques similar to those described above. For example, DR-DOS 7.02 and higher incorporate a safety feature developed by Matthias R. Paul: If these files are called inappropriately, tiny embedded stubs will just display some file version information and exit gracefully. Additionally, the message is specifically crafted to follow certain "magic" patterns recognized by the external NetWare & DR-DOS VERSION file identification utility.A similar protection feature was the 8080 instruction C7h ("RST 0") at the very start of Jay Sage's and Joe Wright's Z-System type-3 and type-4 "Z3ENV" programs as well as "Z3TXT" language overlay files, which would result in a warm boot (instead of a crash) under CP/M-80 if loaded inappropriately.In a distantly similar fashion, many (binary) file formats by convention include a 1Ah byte (ASCII ^Z) near the beginning of the file. This control character will be interpreted as "soft" end-of-file (EOF) marker when a file is opened in non-binary mode, and thus, under many operating systems (including the PDP-6 monitor and RT-11, VMS, TOPS-10, CP/M, DOS, and Windows), it prevents "binary garbage" from being displayed when a file is accidentally printed at the console.
Linux:
FatELF: Universal binaries for Linux FatELF was a fat binary implementation for Linux and other Unix-like operating systems. Technically, a FatELF binary was a concatenation of ELF binaries with some meta data indicating which binary to use on what architecture. Additionally to the CPU architecture abstraction (byte order, word size, CPU instruction set, etc.), there is the advantage of binaries with support for multiple kernel ABIs and versions.
Linux:
FatELF has several use-cases, according to developers: Distributions no longer need to have separate downloads for various platforms.
Separated /lib, /lib32 and /lib64 trees are not required anymore in OS directory structure.
The correct binary and libraries are centrally chosen by the system instead of shell scripts.
If the ELF ABI changes someday, legacy users can be still supported.
Distribution of web browser plug ins that work out of the box with multiple platforms.
Distribution of one application file that works across Linux and BSD OS variants, without a platform compatibility layer on them.
One hard drive partition can be booted on different machines with different CPU architectures, for development and experimentation. Same root file system, different kernel and CPU architecture.
Applications provided by network share or USB sticks, will work on multiple systems. This is also helpful for creating portable applications and also cloud computing images for heterogeneous systems.A proof-of-concept Ubuntu 9.04 image is available. As of 2021, FatELF has not been integrated into the mainline Linux kernel.
Windows:
Fatpack Although the Portable Executable format used by Windows does not allow assigning code to platforms, it is still possible to make a loader program that dispatches based on architecture. This is because desktop versions of Windows on ARM have support for 32-bit x86 emulation, making it a useful "universal" machine code target. Fatpack is a loader that demonstrates the concept: it includes a 32-bit x86 program that tries to run the executables packed into its resource sections one by one.
Windows:
Arm64X When developing Windows 11 ARM64, Microsoft introduced a new way to extend the Portable Executable format called Arm64X. An Arm64X binary contains all the content that would be in separate x64/Arm64EC and Arm64 binaries, but merged into one more efficient file on disk. Visual C++ toolset has been upgraded to support producing such binaries. And when building Arm64X binaries are technically difficult, developers can build Arm64X pure forwarder DLLs instead.
Similar concepts:
The following approaches are similar to fat binaries in that multiple versions of machine code of the same purpose are provided in the same file.
Similar concepts:
Heterogeneous computing Since 2007, some specialized compilers for heterogeneous platforms produce code files for parallel execution on multiple types of processors, i.e. the CHI (C for Heterogeneous Integration) compiler from the Intel EXOCHI (Exoskeleton Sequencer) development suite extends the OpenMP pragma concept for multithreading to produce fat binaries containing code sections for different instruction set architectures (ISAs) from which the runtime loader can dynamically initiate the parallel execution on multiple available CPU and GPU cores in a heterogeneous system environment.Introduced in February 2007, Nvidia's parallel computing platform CUDA (Compute Unified Device Architecture) is a software to enable general-purpose computing on GPUs (GPGPU). Its LLVM-based compiler NVCC can create ELF-based fat binaries containing so called PTX virtual assembly (as text) which the CUDA runtime driver can later just-in-time compile into some SASS (Streaming Assembler) binary executable code for the actually present target GPU. The executables can also include so called CUDA binaries (aka cubin files) containing dedicated executable code sections for one or more specific GPU architectures from which the CUDA runtime can choose from at load-time. Fat binaries are also supported by GPGPU-Sim, a GPU simulator introduced in 2007 as well.Multi2Sim (M2S), an OpenCL heterogeneous system simulator framework (originally only for either MIPS or x86 CPUs, but later extended to also support ARM CPUs and GPUs like the AMD/ATI Evergreen & Southern Islands as well as Nvidia Fermi & Kepler families) supports ELF-based fat binaries as well.
Similar concepts:
Fat objects GNU Compiler Collection (GCC) and LLVM do not have a fat binary format, but they do have fat object files for link-time optimization (LTO). Since LTO involves delaying the compilation to link-time, the object files must store the intermediate representation (IR), but on the other hand machine code may need to be stored too (for speed or compatibility). An LTO object containing both IR and machine code is known as a fat object.
Similar concepts:
Function multi-versioning Even in a program or library intended for the same instruction set architecture, a programmer may wish to make use of some newer instruction set extensions while keeping compatibility with an older CPU. This can be achieved with function multi-versioning (FMV): versions of the same function are written into the program, and a piece of code decides which one to use by detecting the CPU's capabilities (such as through CPUID). Intel C++ Compiler, GCC, and LLVM all have the ability to automatically generate multi-versioned functions. This is a form of dynamic dispatch without any semantic effects.
Similar concepts:
Many math libraries feature hand-written assembly routines that are automatically chosen according to CPU capability. Examples include glibc, Intel MKL, and OpenBLAS. In addition, the library loader in glibc supports loading from alternative paths for specific CPU features.A similar, but byte-level granular approach originally devised by Matthias R. Paul and Axel C. Frinke is to let a small self-discarding, relaxing and relocating loader embedded into the executable file alongside any number of alternative binary code snippets conditionally build a size- or speed-optimized runtime image of a program or driver necessary to perform (or not perform) a particular function in a particular target environment at load-time through a form of dynamic dead code elimination (DDCE). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Dimond ring**
Dimond ring:
A Dimond ring or Dimond ring translator was an early type of computer memory, created in the early 1940s by T. L. Dimond at Bell Laboratories for Bell's #5 Crossbar Switch, a type of early telephone switch.
Structure:
Large-diameter magnetic ferrite toroidal rings with solenoid windings, through which are threaded writing and reading wires.
Uses:
It was used in the #5 Crossbar Switch and TXE (prior to version 4) telephone exchanges. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Inferior hypogastric plexus**
Inferior hypogastric plexus:
The inferior hypogastric plexus (pelvic plexus in some texts) is a network (plexus) of nerves that supplies the organs of the pelvic cavity. The inferior hypogastric plexus gives rise to the prostatic plexus in males and the uterovaginal plexus in females.The inferior hypogastric plexus is a paired structure, meaning there is one on the left and the right side of the body. These are located on either side of the rectum in males, and at the sides of the rectum and vagina in females. For this reason, injury to this structure can arise as a complication of pelvic surgeries and may cause urinary dysfunction and urinary incontinence. Testing of bladder function is used in that case to show a poorly compliant bladder, with bladder neck incompetence, and fixed external sphincter tone.
Structure:
The plexus is formed from: a continuation of the superior hypogastric plexus on either side, at the sacral promontory in the interiliac triangle. At this location, the presacral nerve sits in the middle in only 25% of people and is more commonly present on the left.
sacral splanchnic nerves, from the sympathetic trunk.
pelvic splanchnic nerves (from the second, third, and fourth sacral nerves) also contribute parasympathetic efferent fibers to the plexus.From these plexuses numerous branches are distributed to the viscera of the pelvis.
They accompany the branches of the internal iliac artery.
It is the source for the middle rectal plexus, vesical plexus, prostatic plexus, and uterovaginal plexus. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Nucleoporin 35**
Nucleoporin 35:
Nucleoporin 35 (Nup35) is a protein that in humans is encoded by the NUP35 gene.
Background:
This gene encodes a member of the nucleoporin family. The protein is localized to the nuclear rim and is part of the nuclear pore complex (NPC). All molecules entering or leaving the nucleus either diffuse through or are actively transported by the NPC. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Doubly connected edge list**
Doubly connected edge list:
The doubly connected edge list (DCEL), also known as half-edge data structure, is a data structure to represent an embedding of a planar graph in the plane, and polytopes in 3D. This data structure provides efficient manipulation of the topological information associated with the objects in question (vertices, edges, faces). It is used in many algorithms of computational geometry to handle polygonal subdivisions of the plane, commonly called planar straight-line graphs (PSLG). For example, a Voronoi diagram is commonly represented by a DCEL inside a bounding box. This data structure was originally suggested by Muller and Preparata for representations of 3D convex polyhedra. Later, a somewhat different data structure was suggested, but the name "DCEL" was retained.
Doubly connected edge list:
For simplicity, only connected graphs are considered, however the DCEL structure may be extended to handle disconnected graphs as well by introducing dummy edges between disconnected components.
Data structure:
DCEL is more than just a doubly linked list of edges. In the general case, a DCEL contains a record for each edge, vertex and face of the subdivision. Each record may contain additional information, for example, a face may contain the name of the area. Each edge usually bounds two faces and it is, therefore, convenient to regard each edge as two "half-edges" (which are represented by the two edges with opposite directions, between two vertices, in the picture on the right). Each half-edge is "associated" with a single face and thus has a pointer to that face. All half-edges associated with a face are clockwise or counter-clockwise. For example, in the picture on the right, all half-edges associated with the middle face (i.e. the "internal" half-edges) are counter-clockwise. A half-edge has a pointer to the next half-edge and previous half-edge of the same face. To reach the other face, we can go to the twin of the half-edge and then traverse the other face. Each half-edge also has a pointer to its origin vertex (the destination vertex can be obtained by querying the origin of its twin, or of the next half-edge).
Data structure:
Each vertex contains the coordinates of the vertex and also stores a pointer to an arbitrary edge that has the vertex as its origin. Each face stores a pointer to some half-edge of its outer boundary (if the face is unbounded then pointer is null). It also has a list of half-edges, one for each hole that may be incident within the face. If the vertices or faces do not hold any interesting information, there is no need to store them, thus saving space and reducing the data structure's complexity. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Bioreactor landfill**
Bioreactor landfill:
Landfills are the primary method of waste disposal in many parts of the world, including United States and Canada. Bioreactor landfills are expected to reduce the amount of and costs associated with management of leachate, to increase the rate of production of methane (natural gas) for commercial purposes and reduce the amount of land required for land-fills. Bioreactor landfills are monitored and manipulate oxygen and moisture levels to increase the rate of decomposition by microbial activity.
Traditional landfills and associated problems:
Landfills are the oldest known method of waste disposal. Waste is buried in large dug out pits (unless naturally occurring locations are available) and covered. Bacteria and archaea decompose the waste over several decades producing several by-products of importance, including methane gas (natural gas), leachate, and volatile organic compounds (such as hydrogen sulfide (H2S), N2O2, etc.).
Traditional landfills and associated problems:
Methane gas, a strong greenhouse gas, can build up inside the landfill leading to an explosion unless released from the cell. Leachate are fluid metabolic products from decomposition and contain various types of toxins and dissolved metallic ions. If leachate escapes into the ground water it can cause health problems in both animals and plants. The volatile organic compounds (VOCs) are associated with causing smog and acid rain. With the increasing amount of waste produced, appropriate places to safely store it have become difficult to find.
Working of a bioreactor landfill:
There are three types of bioreactor: aerobic, anaerobic and a hybrid (using both aerobic and anaerobic method). All three mechanisms involve the reintroduction of collected leachate supplemented with water to maintain moisture levels in the landfill. The micro-organisms responsible for decomposition are thus stimulated to decompose at an increased rate with an attempt to minimise harmful emissions.In aerobic bioreactors air is pumped into the landfill using either vertical or horizontal system of pipes. The aerobic environment decomposition is accelerated and amount of VOCs, toxicity of leachate and methane are minimised. In anaerobic bioreactors with leachate being circulated the landfill produces methane at a rate much faster and earlier than traditional landfills. The high concentration and quantity of methane allows it to be used more efficiently for commercial purposes while reducing the time that the landfill needs to be monitored for methane production. Hybrid bioreactors subject the upper portions of the landfill through aerobic-anaerobic cycles to increase decomposition rate while methane is produced by the lower portions of the landfill. Bioreactor landfills produce lower quantities of VOCs than traditional landfills, except H2S. Bioreactor landfills produce higher quantities of H2S. The exact biochemical pathway responsible for this increase is not well studied
Advantages of bioreactor landfills:
Bioreactor landfills accelerate the process of decomposition. As decomposition progresses, the mass of biodegradable components in the landfill declines, creating more space for dumping garbage. Bioreactor landfills are expected to increase this rate of decomposition and save up to 30% of space needed for landfills. With increasing amounts of solid waste produced every year and scarcity of landfill spaces, bioreactor landfill can thus provide a significant way of maximising landfill space. This is not just cost effective, but since less land is needed for the landfills, this is also better for the environment.Furthermore, most landfills are monitored for at least 3 to 4 decades to ensure that no leachate or landfill gases escape into the community surrounding the landfill site. In contrast, bioreactor landfill are expected to decompose to level that does not require monitoring in less than a decade. Hence, the landfill land can be used for other purposes such as reforestation or parks, depending on the location at an earlier date. In addition, re-using leachate to moisturise the landfill filters it. Thus, less time and energy is required to process the leachate, making the process more efficient.
Disadvantages of bioreactor landfills:
Bioreactor landfills are a relatively new technology. For the newly developed bioreactor landfills initial monitoring costs are higher to ensure that everything important is discovered and properly controlled. This includes gases, odours and seepage of leachate into the ground surface.
The increased moisture content of bioreactor landfill may reduce the structural stability of the landfill by increasing the pore water pressure within the waste mass.Since the target of bioreactor landfills is to maintain a high moisture content, gas collection systems can be affected by the increased moisture content of the waste.
Implementation of bioreactor landfills:
Bioreactor landfills being a novel technology are still in the development phase and are being studied in the laboratory-scale. Pilot projects for bioreactor landfills are showing promise and more are being experimented with in different parts of the world. Despite the potential benefits of bioreactor landfills there are no standardised and approved designs with guidelines and operational procedures. Following is a list of bioreactor landfill projects which are being used to collect data for forming these needed guidelines and procedures: United States California Yolo County Florida Alachua County Southeast Landfill Highlands County New River Regional Landfill, Raiford Polk County Landfill, Winter Haven Kentucky Outer Loop Landfill Michigan Saint Clair County Mississippi Plantation Oaks Bioreactor Demonstration Project, Sibley Missouri Columbia New Jersey ACUA's Haneman Environmental Park, Egg Harbor Township North Carolina Buncombe County Landfill Project Virginia Maplewood Landfill and King George County Landfills Virginia Landfill Project XL Demonstration Project Canada Sainte-Sophie Bioreactor demonstration Project, Quebec Australia New South Wales WoodLawn, Goulburn Queensland Ti Tree Bioenergy, Ipswich | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Roland MC-303**
Roland MC-303:
The Roland MC-303 is the first of a series of musical instruments known as a groovebox. It combines a simple sound module with a sequencer to record and store notation, along with controls aimed at encouraging the musician to improvise the music while it is playing. Despite the number in its name and the attention it received at its launch, the MC-303 has more in common with other MC prefixed synthesizers (such as the Roland MC-202), which contain built-in sequencers, than it does with the famous Roland TB-303. As the first Groovebox, the MC-303 was the first in a line of inexpensive products specifically targeted towards house DJs and amateur home musicians rather than professional producers. It was superseded by the Roland MC-505. It is the predecessor to the Roland JX-305, Roland D2, Roland MC-307, Roland EG-101, Roland MC-09, Roland MC-909, Roland MC-808, and most recently the Roland MC-707 in 2019, along with its more portable sibling, the Roland MC-101.
Features:
The key features of the MC-303 are: Sound generator with 28 note - voice polyphony based on the structure model of Roland JV-80 synthesizer 8-track sequencer containing multiple quantize functions: Grid, Shuffle, and Groove. (7 pitched instruments and 1 drum kit) 16-part multitimbral 448 preset sounds and 12 drum kits (includes the Roland CR-78, TR-808, TR-606 & TR-909, electro, techno, jungle, house, drum & bass, breakbeat), 40 synth basses (TB-303, etc.), 35 synth leads, 33 synth pads Resonant filter, LFO, envelope control and built-in effects: delay, reverb, flanger and chorus Realtime Phrase Sequencer (RPS) for instant recall of musical phrases Low Boost Knob feature (Back panel, Only on the Roland MC-303): allowing you to dial in as much low-end as it takes to create powerful Kick or TR-808 Bass sounds so that anyone can 'feel' the groove.
Features:
300 onboard dance music variation patterns such as drumbeats and basslines Recording length of up to 32 bars per pattern Instant storage of up to 50 user patterns, 300 pattern variations and 10 songs Storage space for up to approximately 14,000 notes MIDI in and out connections (but no MIDI thru)
Synthesizer/Sound Module:
The synthesizer built into the Roland MC-303 is a rompler which contains sounds largely drawn from classic Roland synths and drum machines such as the TB-303, TR-808 and TR-909 along with the Juno series and various other dance themed sounds such as pads, pianos, strings and vinyl scratches. The sounds can be manipulated with a low-pass filter, various modulation capabilities and some simple DSP effects. It doesn't have a sampler, although the instruction book contains instructions for getting it to control an external sampler.
Sequencer:
The most important part of the MC-303 is its built-in pattern based 8 track sequencer. Each pattern can contain up to 32 bars. It can record and send MIDI data via the MIDI jacks on the rear panel, enabling its internal sequencer to control other sound modules, or its internal sound module to be controlled by an external sequencer. Although communication with other devices is possible, the main advantage to the MC-303 with its small form factor and all-in-one design is the ability to use it as a self-contained studio, albeit an amateur one. Featuring a micro-keyboard that can also be used as a drum sequencer, the MC-303 imitates the handling as well as the look and feel of other famous Roland synthesizers and drum machines such as the MC-202, TB-303, TR-808 and TR-909.
Criticisms:
One criticism made of the machine in various reviews, including the August 1996 issue of Sound on Sound magazine, was that the sound module was essentially limited to only playing built-in preset sounds, discouraging innovation. From a more technical perspective, a major concern was that any knob tweaks made during real time recording were not transmitted via MIDI. The number of preset patterns (mostly aimed for Trance and Techno music) outweighed the number of programmable user patterns which also discouraged innovation. Roland responded in part to these criticisms in its later grooveboxes by solving the MIDI problem, increasing the synthesis capabilities and user pattern storage and adding a sampler section. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Jonathan S. Turner**
Jonathan S. Turner:
Jonathan Shields Turner is a senior professor of Computer Science in the School of Engineering and Applied Science at Washington University in St. Louis. His research interests include the design and analysis of high performance routers and switching systems, extensible communication networks via overlay networks, and probabilistic performance of heuristic algorithms for NP-complete problems.
Biography:
Jonathan Shields Turner was born on November 13, 1953, in Boston. Turner started his undergraduate studies at Oberlin College, and later enrolled in the undergraduate engineering program at Washington University. In doing so, he became one of the first dual-degree engineering graduates from Washington University. In 1975, he graduated with a B.A. in Theater from Oberlin College. Then, in 1977, he graduated with a B.S. in Computer Science and B.S. in Electrical Engineering from Washington UniversityOnce Turner graduated, he began attending Northwestern University for Computer Science graduate school, and simultaneously began working at Bell Labs as a member of their technical staff. In 1979, he received his M.S. in computer science from Northwestern, and continued on as a doctoral student under the supervision of Hal Sudborough. From 1981 to 1983, he became the principal system architect for the Fast Packet Switching project at Bell Labs. He received eleven patents for his work on the Fast Packet Switching project. In 1982 he published his doctoral dissertation, receiving his Ph.D. in computer science from Northwestern.
Biography:
Turner joined Washington University in 1983 as an assistant professor in the Computer Science and Electrical Engineering departments. In 1986, he published a paper titled "New Directions in Communications (or Which Way to the Information Age)", which forecast the convergence of data, voice, and video traffic on networks, and proposed scalable switching architectures to handle such a traffic load. This paper would later be reprinted in the 50th anniversary issue of the IEEE Communications Magazine as a "landmark article". In 1988 he founded the Advanced Networking Group and co-founded the Applied Research Laboratory with Washington University colleagues Jerome R. Cox and Guru Parulkar. Turner directed the Applied Research Laboratory (ARL) from its inception to 2012, and directed the Advanced Networking group until it was subsumed by the ARL in 1992. He was promoted to full professor by 1990. He became the Computer Science department chair in 1992 and held this position through 1997. In 1998 Turner co-founded a company named Growth Networks—again in collaboration with Professors Jerome Cox and Guru Parulkar—which focused on high performance switching components for Internet routers and Asynchronous Transfer Mode switches. Turner was Chief Scientist at Growth Networks. In 2000 Cisco acquired Growth Networks for $355 million in stock, largely for the intellectual property and engineering talent. At the time of acquisition, Growth Networks had 55 employees. From 2007 to 2008 he again served as department chair of the Computer Science department. Turner retired from Washington University in 2014 after 30 years with the department. He is now a Senior Professor for the department, and still likes to perform research when he is not sailing the Florida coast or playing tennis with his wife.
Awards and distinctions:
Jonathan S. Turner has been awarded 30 patents for his work in switching systems, and has many widely cited publications.Turner has received honors from a variety of professional organizations. In 1990 he was elected as an IEEE Fellow for "contributions to multipoint switching networks for high-speed packetized information transmission". In 1994 he received the IEEE Koji Kobayashi Computers and Communications Award for "fundamental contributions to communications and computing through architectural innovation in high-speed packet networks. In 2000 he was awarded the IEEE Millennium Medal In 2001 he was elected as an ACM Fellow for research involving and extending his 1986 seminal paper. In 2002 he was awarded the James B. Eads Award from the St. Louis Academy of Science, for outstanding achievement in engineering or technology. In 2007 he was elected to the National Academy of Engineering.
Awards and distinctions:
Turner has also received many honors from Washington University. In 1993 he was honored with the Founder's Day Distinguished Faculty Award, which is awarded to faculty who have an "outstanding commitment to the intellectual and personal development of students". In 1994 he became the Henry Edwin Sever Chair of Engineering, which at that time was a new endowed professorship. He held this position until 2006. In 2004 he won the Arthur Holly Compton Faculty Achievement Award, which is similar to the Founder's Day Distinguished Faculty Award but more selective. In 2006 Turner was named the Barbara J. and Jerome R. Cox Professor of Computer Science for "advancing the relationship between theory and practice in the design of digital systems." In 2007 he received an Alumni Achievement Award from the School of Engineering and Applied Science. In 2014 he received the Dean's Award from the Dean of the School of Engineering and Applied Science. Also that year the Computer Science department created the Turner Dissertation Award in recognition of his many achievements and research contributions. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Supergravity**
Supergravity:
In theoretical physics, supergravity (supergravity theory; SUGRA for short) is a modern field theory that combines the principles of supersymmetry and general relativity; this is in contrast to non-gravitational supersymmetric theories such as the Minimal Supersymmetric Standard Model. Supergravity is the gauge theory of local supersymmetry. Since the supersymmetry (SUSY) generators form together with the Poincaré algebra a superalgebra, called the super-Poincaré algebra, supersymmetry as a gauge theory makes gravity arise in a natural way.
Gravitons:
Like all covariant approaches to quantum gravity, supergravity contains a spin-2 field whose quantum is the graviton. Supersymmetry requires the graviton field to have a superpartner. This field has spin 3/2 and its quantum is the gravitino. The number of gravitino fields is equal to the number of supersymmetries.
History:
Gauge supersymmetry The first theory of local supersymmetry was proposed by Dick Arnowitt and Pran Nath in 1975 and was called gauge supersymmetry.
History:
Supergravity The first model of 4-dimensional supergravity (without this denotation) was formulated by Dmitri Vasilievich Volkov and Vyacheslav A. Soroka in 1973, emphasizing the importance of spontaneous supersymmetry breaking for the possibility of a realistic model. The minimal version of 4-dimensional supergravity (with unbroken local supersymmetry) was constructed in detail in 1976 by Dan Freedman, Sergio Ferrara and Peter van Nieuwenhuizen. In 2019 the three were awarded a special Breakthrough Prize in Fundamental Physics for the discovery. The key issue of whether or not the spin 3/2 field is consistently coupled was resolved in the nearly simultaneous paper, by Deser and Zumino, which independently proposed the minimal 4-dimensional model. It was quickly generalized to many different theories in various numbers of dimensions and involving additional (N) supersymmetries. Supergravity theories with N>1 are usually referred to as extended supergravity (SUEGRA). Some supergravity theories were shown to be related to certain higher-dimensional supergravity theories via dimensional reduction (e.g. N=1, 11-dimensional supergravity is dimensionally reduced on T7 to 4-dimensional, ungauged, N = 8 supergravity). The resulting theories were sometimes referred to as Kaluza–Klein theories as Kaluza and Klein constructed in 1919 a 5-dimensional gravitational theory, that when dimensionally reduced on a circle, its 4-dimensional non-massive modes describe electromagnetism coupled to gravity.
History:
mSUGRA mSUGRA means minimal SUper GRAvity. The construction of a realistic model of particle interactions within the N = 1 supergravity framework where supersymmetry (SUSY) breaks by a super Higgs mechanism carried out by Ali Chamseddine, Richard Arnowitt and Pran Nath in 1982. Collectively now known as minimal supergravity Grand Unification Theories (mSUGRA GUT), gravity mediates the breaking of SUSY through the existence of a hidden sector. mSUGRA naturally generates the Soft SUSY breaking terms which are a consequence of the Super Higgs effect. Radiative breaking of electroweak symmetry through Renormalization Group Equations (RGEs) follows as an immediate consequence. Due to its predictive power, requiring only four input parameters and a sign to determine the low energy phenomenology from the scale of Grand Unification, its interest is a widely investigated model of particle physics 11D: the maximal SUGRA One of these supergravities, the 11-dimensional theory, generated considerable excitement as the first potential candidate for the theory of everything. This excitement was built on four pillars, two of which have now been largely discredited: Werner Nahm showed 11 dimensions as the largest number of dimensions consistent with a single graviton, and more dimensions will show particles with spins greater than 2. However, if two of these dimensions are time-like, these problems are avoided in 12 dimensions. Itzhak Bars gives this emphasis.
History:
In 1981 Ed Witten showed 11 as the smallest number of dimensions big enough to contain the gauge groups of the Standard Model, namely SU(3) for the strong interactions and SU(2) times U(1) for the electroweak interactions. Many techniques exist to embed the standard model gauge group in supergravity in any number of dimensions like the obligatory gauge symmetry in type I and heterotic string theories, and obtained in type II string theory by compactification on certain Calabi–Yau manifolds. The D-branes engineer gauge symmetries too.
History:
In 1978 Eugène Cremmer, Bernard Julia and Joël Scherk (CJS) found the classical action for an 11-dimensional supergravity theory. This remains today the only known classical 11-dimensional theory with local supersymmetry and no fields of spin higher than two. Other 11-dimensional theories known and quantum-mechanically inequivalent reduce to the CJS theory when one imposes the classical equations of motion. However, in the mid 1980s Bernard de Wit and Hermann Nicolai found an alternate theory in D=11 Supergravity with Local SU(8) Invariance. While not manifestly Lorentz-invariant, it is in many ways superior, because it dimensionally-reduces to the 4-dimensional theory without recourse to the classical equations of motion.In 1980 Peter Freund and M. A. Rubin showed that compactification from 11 dimensions preserving all the SUSY generators could occur in two ways, leaving only 4 or 7 macroscopic dimensions, the others compact. The noncompact dimensions have to form an anti-de Sitter space. There are many possible compactifications, but the Freund-Rubin compactification's invariance under all of the supersymmetry transformations preserves the action.Finally, the first two results each appeared to establish 11 dimensions, the third result appeared to specify the theory, and the last result explained why the observed universe appears to be four-dimensional.
History:
Many of the details of the theory were fleshed out by Peter van Nieuwenhuizen, Sergio Ferrara and Daniel Z. Freedman.
History:
The end of the SUGRA era The initial excitement over 11-dimensional supergravity soon waned, as various failings were discovered, and attempts to repair the model failed as well. Problems included: The compact manifolds which were known at the time and which contained the standard model were not compatible with supersymmetry, and could not hold quarks or leptons. One suggestion was to replace the compact dimensions with the 7-sphere, with the symmetry group SO(8), or the squashed 7-sphere, with symmetry group SO(5) times SU(2).
History:
Until recently, the physical neutrinos seen in experiments were believed to be massless, and appeared to be left-handed, a phenomenon referred to as the chirality of the Standard Model. It was very difficult to construct a chiral fermion from a compactification — the compactified manifold needed to have singularities, but physics near singularities did not begin to be understood until the advent of orbifold conformal field theories in the late 1980s.
History:
Supergravity models generically result in an unrealistically large cosmological constant in four dimensions, and that constant is difficult to remove, and so require fine-tuning. This is still a problem today.
History:
Quantization of the theory led to quantum field theory gauge anomalies rendering the theory inconsistent. In the intervening years physicists have learned how to cancel these anomalies.Some of these difficulties could be avoided by moving to a 10-dimensional theory involving superstrings. However, by moving to 10 dimensions one loses the sense of uniqueness of the 11-dimensional theory.The core breakthrough for the 10-dimensional theory, known as the first superstring revolution, was a demonstration by Michael B. Green, John H. Schwarz and David Gross that there are only three supergravity models in 10 dimensions which have gauge symmetries and in which all of the gauge and gravitational anomalies cancel. These were theories built on the groups SO(32) and E8×E8 , the direct product of two copies of E8. Today we know that, using D-branes for example, gauge symmetries can be introduced in other 10-dimensional theories as well.
History:
The second superstring revolution Initial excitement about the 10-dimensional theories, and the string theories that provide their quantum completion, died by the end of the 1980s. There were too many Calabi–Yaus to compactify on, many more than Yau had estimated, as he admitted in December 2005 at the 23rd International Solvay Conference in Physics. None quite gave the standard model, but it seemed as though one could get close with enough effort in many distinct ways. Plus no one understood the theory beyond the regime of applicability of string perturbation theory.
History:
There was a comparatively quiet period at the beginning of the 1990s; however, several important tools were developed. For example, it became apparent that the various superstring theories were related by "string dualities", some of which relate weak string-coupling - perturbative - physics in one model with strong string-coupling - non-perturbative - in another.
Then the second superstring revolution occurred. Joseph Polchinski realized that obscure string theory objects, called D-branes, which he discovered six years earlier, equate to stringy versions of the p-branes known in supergravity theories. String theory perturbation didn't restrict these p-branes. Thanks to supersymmetry, p-branes in supergravity gained understanding well beyond the limits of string theory.
History:
Armed with this new nonperturbative tool, Edward Witten and many others could show all of the perturbative string theories as descriptions of different states in a single theory that Witten named M-theory. Furthermore, he argued that M-theory's long wavelength limit, i.e. when the quantum wavelength associated to objects in the theory appear much larger than the size of the 11th dimension, needs 11-dimensional supergravity descriptors that fell out of favor with the first superstring revolution 10 years earlier, accompanied by the 2- and 5-branes.
History:
Therefore, supergravity comes full circle and uses a common framework in understanding features of string theories, M-theory, and their compactifications to lower spacetime dimensions.
Relation to superstrings:
The term "low energy limits" labels some 10-dimensional supergravity theories. These arise as the massless, tree-level approximation of string theories. True effective field theories of string theories, rather than truncations, are rarely available. Due to string dualities, the conjectured 11-dimensional M-theory is required to have 11-dimensional supergravity as a "low energy limit". However, this doesn't necessarily mean that string theory/M-theory is the only possible UV completion of supergravity; supergravity research is useful independent of those relations.
4D N = 1 SUGRA:
Before we move on to SUGRA proper, let's recapitulate some important details about general relativity. We have a 4D differentiable manifold M with a Spin(3,1) principal bundle over it. This principal bundle represents the local Lorentz symmetry. In addition, we have a vector bundle T over the manifold with the fiber having four real dimensions and transforming as a vector under Spin(3,1).
4D N = 1 SUGRA:
We have an invertible linear map from the tangent bundle TM to T. This map is the vierbein. The local Lorentz symmetry has a gauge connection associated with it, the spin connection.
4D N = 1 SUGRA:
The following discussion will be in superspace notation, as opposed to the component notation, which isn't manifestly covariant under SUSY. There are actually many different versions of SUGRA out there which are inequivalent in the sense that their actions and constraints upon the torsion tensor are different, but ultimately equivalent in that we can always perform a field redefinition of the supervierbeins and spin connection to get from one version to another.
4D N = 1 SUGRA:
In 4D N=1 SUGRA, we have a 4|4 real differentiable supermanifold M, i.e. we have 4 real bosonic dimensions and 4 real fermionic dimensions. As in the nonsupersymmetric case, we have a Spin(3,1) principal bundle over M. We have an R4|4 vector bundle T over M. The fiber of T transforms under the local Lorentz group as follows; the four real bosonic dimensions transform as a vector and the four real fermionic dimensions transform as a Majorana spinor. This Majorana spinor can be reexpressed as a complex left-handed Weyl spinor and its complex conjugate right-handed Weyl spinor (they're not independent of each other). We also have a spin connection as before.
4D N = 1 SUGRA:
We will use the following conventions; the spatial (both bosonic and fermionic) indices will be indicated by M, N, ... . The bosonic spatial indices will be indicated by μ, ν, ..., the left-handed Weyl spatial indices by α, β,..., and the right-handed Weyl spatial indices by α˙ , β˙ , ... . The indices for the fiber of T will follow a similar notation, except that they will be hatted like this: M^,α^ . See van der Waerden notation for more details. M=(μ,α,α˙) . The supervierbein is denoted by eNM^ , and the spin connection by ωM^N^P . The inverse supervierbein is denoted by EM^N The supervierbein and spin connection are real in the sense that they satisfy the reality conditions eNM^(x,θ¯,θ)∗=eN∗M^∗(x,θ,θ¯) where μ∗=μ , α∗=α˙ , and α˙∗=α and ω(x,θ¯,θ)∗=ω(x,θ,θ¯) .The covariant derivative is defined as DM^f=EM^N(∂Nf+ωN[f]) .The covariant exterior derivative as defined over supermanifolds needs to be super graded. This means that every time we interchange two fermionic indices, we pick up a +1 sign factor, instead of -1.
4D N = 1 SUGRA:
The presence or absence of R symmetries is optional, but if R-symmetry exists, the integrand over the full superspace has to have an R-charge of 0 and the integrand over chiral superspace has to have an R-charge of 2.
A chiral superfield X is a superfield which satisfies D¯α˙^X=0 . In order for this constraint to be consistent, we require the integrability conditions that {D¯α˙^,D¯β˙^}=cα˙^β˙^γ˙^D¯γ˙^ for some coefficients c.
4D N = 1 SUGRA:
Unlike nonSUSY GR, the torsion has to be nonzero, at least with respect to the fermionic directions. Already, even in flat superspace, Dα^eα˙^+D¯α˙^eα^≠0 In one version of SUGRA (but certainly not the only one), we have the following constraints upon the torsion tensor: Tα_^β_^γ_^=0 Tα^β^μ^=0 Tα˙^β˙^μ^=0 Tα^β˙^μ^=2iσα^β˙^μ^ Tμ^α_^ν^=0 Tμ^ν^ρ^=0 Here, α_ is a shorthand notation to mean the index runs over either the left or right Weyl spinors.
4D N = 1 SUGRA:
The superdeterminant of the supervierbein, |e| , gives us the volume factor for M. Equivalently, we have the volume 4|4-superform eμ^=0∧⋯∧eμ^=3∧eα^=1∧eα^=2∧eα˙^=1∧eα˙^=2 If we complexify the superdiffeomorphisms, there is a gauge where Eα˙^μ=0 , Eα˙^β=0 and Eα˙^β˙=δα˙β˙ . The resulting chiral superspace has the coordinates x and Θ.
R is a scalar valued chiral superfield derivable from the supervielbeins and spin connection. If f is any superfield, (D¯2−8R)f is always a chiral superfield.
The action for a SUGRA theory with chiral superfields X, is given by S=∫d4xd2Θ2E[38(D¯2−8R)e−K(X¯,X)/3+W(X)]+c.c.
where K is the Kähler potential and W is the superpotential, and E is the chiral volume factor.
4D N = 1 SUGRA:
Unlike the case for flat superspace, adding a constant to either the Kähler or superpotential is now physical. A constant shift to the Kähler potential changes the effective Planck constant, while a constant shift to the superpotential changes the effective cosmological constant. As the effective Planck constant now depends upon the value of the chiral superfield X, we need to rescale the supervierbeins (a field redefinition) to get a constant Planck constant. This is called the Einstein frame.
N = 8 supergravity in 4 dimensions:
N = 8 supergravity is the most symmetric quantum field theory which involves gravity and a finite number of fields. It can be found from a dimensional reduction of 11D supergravity by making the size of 7 of the dimensions go to zero. It has 8 supersymmetries which is the most any gravitational theory can have since there are 8 half-steps between spin 2 and spin −2. (A graviton has the highest spin in this theory which is a spin 2 particle). More supersymmetries would mean the particles would have superpartners with spins higher than 2. The only theories with spins higher than 2 which are consistent involve an infinite number of particles (such as string theory and higher-spin theories). Stephen Hawking in his A Brief History of Time speculated that this theory could be the Theory of Everything. However, in later years this was abandoned in favour of string theory. There has been renewed interest in the 21st century with the possibility that this theory may be finite.
Higher-dimensional SUGRA:
Higher-dimensional SUGRA is the higher-dimensional, supersymmetric generalization of general relativity. Supergravity can be formulated in any number of dimensions up to eleven. Higher-dimensional SUGRA focuses upon supergravity in greater than four dimensions.
Higher-dimensional SUGRA:
The number of supercharges in a spinor depends on the dimension and the signature of spacetime. The supercharges occur in spinors. Thus the limit on the number of supercharges cannot be satisfied in a spacetime of arbitrary dimension. Some theoretical examples in which this is satisfied are: 12-dimensional two-time theory 11-dimensional maximal SUGRA 10-dimensional SUGRA theories Type IIA SUGRA: N = (1, 1) IIA SUGRA from 11d SUGRA Type IIB SUGRA: N = (2, 0) Type I gauged SUGRA: N = (1, 0) 9d SUGRA theories Maximal 9d SUGRA from 10d T-duality N = 1 Gauged SUGRAThe supergravity theories that have attracted the most interest contain no spins higher than two. This means, in particular, that they do not contain any fields that transform as symmetric tensors of rank higher than two under Lorentz transformations. The consistency of interacting higher spin field theories is, however, presently a field of very active interest. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Space sustainability**
Space sustainability:
Space sustainability aims to maintain the safety and health of the space environment.Similar to sustainability initiatives on Earth, space sustainability seeks to use the environment of space to meet the current needs of society without compromising the needs of future generations. It usually focuses on space closest to Earth, Low Earth Orbit (LEO), since this environment is the one most used and therefore most relevant to humans. It also considers Geostationary Equatorial Orbit (GEO) as this orbit is another popular choice for Earth-orbiting mission designs.The issue of space sustainability is a new phenomenon that is gaining more attention in recent years as the launching of satellites and other space objects has increased. These launches have resulted in more space debris orbiting Earth, hindering the ability of nations to operate in the space environment while increasing the risk of a future launch-related accident that could disrupt its proper use. Space weather also acts as an outstanding factor for spacecraft failure. The current protocol for spacecraft disposal at end-of-life has, at large, not been followed in mission designs and demands extraneous amounts of time for disposal.Precedent created through prior policy initiatives has facilitated initial mitigation of space pollution and created a foundation for space sustainability efforts. To further mitigation, international and transdisciplinary consortia have stepped forward to analyze existing operations, develop standards, and incentivize future procedures to prioritize a sustainable approach. A shift towards sustainable interactions with the space environment is growing in urgency due to the implications of climate change and increasing risk to spacecraft as time presses on.
Fundamentals:
Space sustainability requires all space participants to have three consensuses. The space field should be used peacefully, jointly protect the space field from harm, and maximize space utilization through environmental, economic, and security exploration of space. These consensuses also clarify the relationship between space sustainability and international security, that states and individuals explore space for various purposes. Their reliance on space needs to be guided by rules, order, and policies and obtain more benefits without negatively affecting the space environment and space activities.However, striking an agreement remains challenging even with such demands in place. In the discussions between countries on long-term sustainability, technical improvements are given more importance than introducing and applying new legal regimes. Specifically, technical approaches to space debris have been proposed, such as debris removal. Specific data on space debris is also being explored to help study its impact on sustainability and promote further cooperation between countries.
Current state:
Space sustainability comes into play to address the pressing current state of near-Earth orbits and its high amounts of orbital debris. Spacecraft collisions with orbital debris, space weather, overcrowding in low earth orbit (LEO) makes spacecraft susceptible to higher rates of failure. The current end-of-life protocol for spacecraft exacerbates the space sustainability crisis; many spacecraft are not properly disposed, which increasing the likelihood of further collisions.
Current state:
Orbital debris Orbital debris is defined as unmanned, inoperate objects that exist in space. This orbital debris breaks down further as time progresses as a result of naturally occurring events, such as high-velocity collisions with micrometeoroids, and forced events, such as a controlled release of a launch vehicle. In LEO, these collisions can take place at speeds anywhere between an average velocity of 9 kilometers per second (km/s) and 14 km/s relative to the debris and spacecraft. In GEO, however, these high-speed collisions are a much lower risk as the average relative velocity between the debris and spacecraft is typically between 0 km/s and 2.5 km/s. As of 2012, the United States Joint Space Operations Center tracked 21,000 pieces of orbital debris larger than 10 cm in Earth's nearby orbits (LEO, GEO, and sun-synchronous), where 16,000 of these pieces are catalogued. Space debris can be categorized into three categories: small, medium, and large. Small debris is for pieces that are less than 10 centimeters (cm). Medium-sized debris is for pieces larger than 10 cm, but not an entire spacecraft. Large-sized debris has no official classification, but typically refers to entire spacecraft, such as an out of use satellite or launch vehicle. It is difficult to track small-sized debris in LEO, and challenging to track small and medium-sized debris in GEO. Yet this statement is not to discount the abilities of LEO and GEO tracking capabilities, the smallest piece of tracked debris can weigh as low as ten grams. If the size of the debris prohibits it from being tracked, it also cannot be avoided by the spacecraft and does not allow the spacecraft to lower its risk of collisions. The likelihood of the Kessler syndrome, which essentially states that each collision produces more debris, grows larger as the amount of orbital debris multiplies, increasing the amount of further collisions until space cannot be used entirely.
Current state:
Space weather Space weather poses a risk to satellite health, consequently, resulting in greater amounts of orbital debris. Space weather impacts satellite health in a variety of ways. Firstly, surface charging from the sun's surface facilitates electrical discharges, damaging on-orbit electronics, posing a threat to mission failure. Single Event Upsets (SEUs) can also damage electronics. Dielectric charging and bulk charging can also occur, causing energy problems within the spacecraft. Additionally, at altitudes less than one thousand kilometers, atmospheric drag can increase during solar storms by increasing the altitude of the spacecraft, only adding more drag onto the spacecraft. These factors degrade performance over the spacecraft's lifetime, leaving the spacecraft more susceptible to further system and mission failures.
Current state:
Overcrowding There has been a dramatic increase in the use of LEO and GEO orbits over the last sixty years since the first satellite launch in 1957. To date, there have been approximately ten thousand satellite launches, whereas only approximately 2000 are still active. These satellites can be used for a variety of purposes, which are telecommunications, navigation, weather monitoring, and exploration. Within the coming decade, companies like SpaceX are predicted to launch an additional fifteen thousand satellites into LEO and GEO orbits. Microsatellites built by universities or research organizations have also increased in popularity, contributing to the overcrowding of near earth orbits. This overcrowding of LEO and GEO orbits increases the likelihood of potential collisions among satellites and orbital debris, contributing further to the large amount of orbital debris present in space.
Current state:
End of life protocol The current end of life protocol is that at the end of mission, spacecraft are either added to the graveyard orbit or at a low enough altitude that drag will allow the spacecraft to burn up upon reentry and fall back to Earth. Approximately twenty satellites are put into the graveyard orbit each year. There is no current process to return satellites to Earth after entering the graveyard orbit. The process of a spacecraft returning to Earth via drag can take between ten and one hundred years. This protocol is critical to reduce overcrowding in near-Earth orbits.
Current state:
Mega constellation and space debris The impact of constellations on the space environment has also been studied, such as the probability of collisions of mega constellations in the presence of large amounts of space debris. Although studies have shown that the predictors of mega constellations are highly variable, specific information related to mega constellations is not transparent.But any catastrophic collision, as in the case of Kessler syndrome, has consequences for people and the environment. Putting this thinking into mega constellations, mega constellations existence may have potential benefits, but it will not bring adequate help to the governance of space debris. At the same time, the space debris situation cannot be underestimated or ignored because of the existence of mega constellations.
Concern:
The existence of orbital debris has caused great trouble to the conduct of space activities. The development of space sustainability has not given sufficient political attention, although some warnings and discussions have made this abundantly clear. Debris management is still voluntary on the part of the state, and there are no laws mandating debris management practices, including the amount of debris to be managed. Although the UN Space Debris Mitigation Guidelines were promulgated in 2007 as an initial measure of space debris governance, there is still no broad consensus or action on further limits on space debris after that.
Concern:
The difficulties for individuals wishing to participate in debris management initiatives cannot be ignored. Any individual or sector desiring to participate in space debris operations needs to obtain permission from the launching state, which is difficult for the launching state to do. This is because the process of space debris management inevitably has a negative impact on other space objects, and there is a lot of subsequent liability in terms of financial consumption. Therefore, the launching state would argue that space debris management requires the joint efforts of all states. However, it is difficult to determine what actions can be taken to gain acceptance between countries.
Regulations:
Current space sustainability efforts rely heavily on the precedent set by regulatory agreements and conventions of the twentieth century. Much of this precedent is included in or is related to the Outer Space Treaty of 1963, which represented one of the initial major efforts by the United Nations to create legal frameworks for the operation of nations in space.
Regulations:
Pre-Outer Space Treaty The international community has had concerns about space contamination since the 1950s prior to the launch of Sputnik I. These concerns stemmed from the idea that increasing rates of exploration into further areas of outer space could lead to contamination capable of damaging other planetary bodies, resulting in limitations to human exploration on these bodies and potential harm to the Earth. Efforts to combat these concerns began in 1956 with the International Astronautical Federation (IAF) and the United Nations Committee on the Peaceful Uses of Outer Space (COPUOUS). These efforts continued to 1957 through the National Academy of Sciences and International Council for Science (ICSU). Each of these organizations aimed to study space contamination and develop strategies for how to best address its potential consequences. The ICSU went on to create the Committee on Contamination by Extraterrestrial Exploration (CETEX) that put forward recommendations leading to the establishment of the Committee on Space Research (COSPAR). COSPAR continues to address outer space research on an international scale today [cite cospar].
Regulations:
Outer Space Treaty Relevant regulations of international space law to sustainability in space can be found in the Outer Space Treaty, which was adopted by the UN General Assembly in 1963. The Outer Space Treaty contains seventeen articles designed to create a basic framework for how international law can be applied in outer space. Basic principles of the Outer Space Treaty include the provision in Article IX that parties should "avoid harmful contamination of space and celestial bodies;" definitions of "harmful contamination" are not provided. Other articles of relevance to space sustainability include articles I, II, and III that concern the fair and inclusive international use of space in a manner free from sovereignty, ownership, or occupation by any nation. In addition, articles VII and VIII protect ownership by their respective countries of any objects launched to space while attributing responsibility for any damages to the property or personnel of other countries by those objects to said countries. Descriptions or definitions for what these damages may entail are not provided.
Regulations:
COSPAR Planetary Protection Policy Principles of Article IX provided the basis for the Committee of Space Research (COSPAR) Planetary Protection Policy guidelines, which are generally well-regarded among scientific experts. Such guidelines, however, are non-binding and often described as "soft-law," as they lack legal mandate. The Planetary Protection Policy is primarily concerned with providing information regarding best practices to avoid contamination of the space environment during space exploration missions. COSPAR believes that the prevention of such contamination is in the best interest of humanity as it may impede scientific progress, exploration, and the mission of a search for life. In addition, the argument is made that cross-contamination of the Earth can be potentially harmful to its environment due to the largely unknown nature of potential space contaminants.
Regulations:
Other relevant regulations Regulatory clarifications concerning the Outer Space Treaty of 1963 of relevance to space sustainability were made in subsequent years. The 1967 Return Agreement relates mainly to the return of lost astronauts to their appropriate nations, but also requires Outer Space Treaty signing nations to assist other nations with the return of objects that return to Earth from orbit to their proper owners The 1972 Liability Convention attributes liability for damages from space objects to the nation that launched the object, regardless of whether that damage occurred in space or on Earth. Other clarifications include the 1975 registration convention that attempted to create mechanisms for nations to identify space objects, and the 1979 Moon Agreement that established protections for the environments of the Moon and other nearby planetary bodies. These agreements and conventions represented attempts to improve the initial Outer Space Treaty as space exploration continued to grow in importance throughout the 20th century.
Attitudes:
Countries and major international institutions Both the state and space agencies are working to improve the laws and regulations that facilitate the long-term sustainability of space. For example, the European Code of Conduct for Space Debris Mitigation signed by France, the UK and other countries in 2016. China, Brazil, Mexico and others have legal background and methodological measures under long-term space sustainability. However, the main problem is that until the concept of space sustainability is agreed between countries, inter-regional efforts are not working well.Currently, the Committee on the Peaceful Uses of Outer Space (COPUOS) encourages states to incorporate the space debris mitigation guidelines developed by bodies such as the Inter-Agency Space Debris Coordination (IADC) into their national legislation, thereby regulating state behavior. Some countries have responded positively to this, such as Switzerland, the Netherlands and Spain. However, there are still some countries that do not consider debris management approaches in their national legislation, such as Japan and Australia. Many delegates at the COPUOS meeting expressed their reasons for doing so, arguing that space debris management is closely linked to technology and funding. Technology is dynamic and constantly evolving. Therefore, the incorporation of debris governance guidelines into national law is not an immediate priority at this time.
Attitudes:
Scientific attitudes A study outlined rationale for governance that regulates the current free externalization of true costs and risks, treating orbital space around the Earth as an "additional ecosystem" or a common "part of the human environment" which should be subject to the same concerns and regulations like oceans on Earth. While scientists may not have the means to make and enforce global laws themselves, the study concluded in 2022 that it needs "new policies, rules and regulations at national and international level".
Mitigation:
Sustainability mitigation efforts include but are not limited to design specifications, policy change, removal of space debris, and restoration of orbiting semi-functional technologies. Efforts begin by regulating the debris released during normal operations and post-mission breakups [6]. Due to the increased awareness of high-velocity collisions and orbital debris in the previous decades, missions have adapted design specifications to account for these risks. For example, the RADARSAT program implemented 17 kilograms of shielding to their spacecraft, which increased the program's predicted success rate to 87% from 50%. Another effort in mitigation is restoring semi-functional satellites, which allows a spacecraft classified as “debris” to “functional.” Space debris mitigation focuses on limiting debris release during normal operations, collisions and intentional destruction. Mitigation also includes reducing the possibility for post-mission breakups due to stored energy and/or operations phases, as well as addressing procedure for end of mission disposal for spacecraft.
Mitigation:
Space Sustainability Rating One example leading the regulatory sustainability measures is the Space Sustainability Rating (SSR), which is an instigator for industry competitors to incorporate sustainability into spacecraft design. The Space Sustainability Rating was first conceptualized at the World Economic Forum Global Future Council on Space Technologies designed by international and transdisciplinary consortia. The four leading organizations are the European Space Agency, Massachusetts Institute of Technology, University of Texas at Austin, and BryceTech with the goal to define the technical and programmatic aspects of the SSR. The SSR represents an innovative approach to combating orbital debris through incentivizing the industry to prioritize sustainable and responsible operations. This response entails the consideration of potential harm to the space environment and other spacecraft, all while maintaining mission objectives and high-quality service. The rating takes inspiration from other standards, like leadership in energy and environmental design (LEED) for the building sector. Several of the factors emphasized in the rating were extracted from LEED design considerations like the incorporation of feedback and public comments, or the rating's advocacy to influence policy, such as orbit fragmentation risks, collision avoidance capabilities, trackability, and adoption of international standards.
Mitigation:
Tracking Tracking is one of the main Space Sustainability Rating modules’ efforts. The module "Detectability, Identification and Tracking" (DIT) consists of standardizing the comparison of satellite missions to encourage satellite operators to improve their satellite design and operational approaches for the observer to detect, identify, and track the satellites. Tracking presents challenges when the observer seeks to monitor and predict the spacecraft behavior over time. While the observer may know the name, owner, and instantaneous location of the satellite, the operator controls the full knowledge of the orbital parameters. The Space Situational Awareness (SSA) is one the tools geared towards solving the challenges presented when tracking orbiting satellites and debris. The SSA continuously tracks objects using ground-based radar and optical stations so the orbital paths of debris can be predicted and operations avoid collisions. It feeds data to 30 different systems like satellites, optical telescopes, radar systems, and supercomputers to predict risk of collision days in advance. Other efforts in tracking orbital debris are made by the US Space Surveillance Network (SSN).
Mitigation:
Removal Under the "External Services" module of the SSR, the rating offers commitment to use or demonstration of use of end-of-life removal services. Space debris mitigation measures are found to be inadequate to stabilize debris environments with an actual current compliance of approximately sixty percent. Moreover, a low compliance rate of approximately thirty percent of the 103 spacecraft that reached end of life between 1997 and 2003 were disposed of in a graveyard orbit. Since policy has not caught up to ensure the longevity of LEO for future generations, actions like Active Debris Removal (ADR) are being considered to stabilize the future of LEO environment. Most famous removal concepts are based on directed energy, momentum exchange or electrodynamics, aerodynamic drag augmentation, solar sails, auxiliary propulsion units, retarding surfaces and on-orbit capture. As ADR consists of an external disposal method to remove obsolete satellites or spacecraft fragments. Since large-sized debris objects in orbit provide a potential source for tens of thousands of fragments in the future, ADR efforts focus on objects with high mass and large cross-sectional areas, in densely populated regions, and at high altitudes; in this instance, retired satellites and rocket bodies are a priority. Other practical advancements toward space debris removal include missions like RemoveDEBRIS and End-of-Life Service (ELS-d).
Growing urgency:
The previous reduced state of regulation and mitigation on space debris and rocket fuel emissions is aggravating the Earth's stratosphere through collisions and ozone depletion, increasing the risk for spacecraft health through its lifetime.
Growing urgency:
Inaccessibility to LEO Due to the increase of satellites being launched and the growing amount of orbital debris in LEO, the risk of LEO becoming inaccessible over time (in accordance with the Kessler syndrome) is increasing in likelihood. The mitigation policies for creating space debris fall under an area of voluntary codes by the states, although it has been disputed whether the Article I Outer Space Treaty or the Article IX Outer Space Treaty protects the space environment from deliberate harm, which has yet to be upheld. In 2007, an inactive Chinese satellite was purposefully destroyed by the Chinese government as a part of their anti-satellite weapon test (ASAT), spreading nearly 2800 objects of space debris five centimeters or larger into LEO. An analysis concluded that about eighty percent of the debris will remain in LEO nine years after this destruction. In addition, the destruction increased the collision likelihood for three Italian satellites that launched the same year as the Fengyun-1C destruction. The increase in collision ranged between ten and sixty percent. However, there were no legal consequences against the Chinese government.
Growing urgency:
Rocket fuel emissions When rockets are launched into space, parts of their fuel enter the stratosphere of the Earth. Rocket fuel emissions are made up of carbon dioxide, water, hydrochloric acid, alumina and soot particles. The most concerning emissions from rocket fuel are chlorine and alumina particles from solid rocket motors (SRMs) and soot from kerosene fueled engines. When the hydrochloric acid from the engine exhaust dissociates, the free chlorine roams freely in the stratosphere. The chemical reaction between these chlorine and alumina causes ozone depletion. In addition, the soot particles form over a black umbrella over the stratosphere which can cause the temperature of the Earth's surface to lower and further depleting the ozone layer, an unintentional form of geoengineering. The nature of geoengineering has been disputed as a form of mitigating global warming and has the possibility of being banned and holding rockets accountable for the soot particles they distribute to the stratosphere. New types of engines and fuels are emerging, mainly the liquid oxygen (LOX) and monomethylhydrazine engine, but there is minimal research on their impact on the environment besides their emission of hydroxide and nitrogen oxide compounds, two molecules that have significant impact on the ozone layer. Currently, rocket fuel emissions have been deemed insignificant when it comes to their consequences to Earth's environment and LEO. However, emissions will increase in the coming years, making rocket fuel's contribution to global warming much more significant.
Beyond LEO:
Space sustainability concepts and mindsets tend to stay in Low Earth Orbit (LEO). One reason that cannot be ignored is that it is easier to discuss the problem at hand than to speculate on the unknown. There are also examples to prove that since Apollo 17 completed its mission and stayed in Low Earth orbit in 1972, human-crewed space missions in Low Earth orbit have ceased to exist. In this way, it is a reasonable assumption that the closer Moon could be the next object to be explored when the gaze is not limited to LEO. Both lunar orbit and LEO are part of the space environment. In the context of the presence of space debris in LEO, it is normal to speculate that lunar orbit also possesses the nuisance of debris. Space debris measures similar to those in LEO related to space sustainability would be taken.Not only has the Moon been the subject of study, but other bodies have also been taken into account. Elon Musk, the chief executive of SpaceX at the International Astronautical Congress in 2016, explained the ambitious goal of exploring Mars in the 22nd century. But complicated issues remain, such as the technical aspects of achieving long-distance space flight and the rules and legal aspects associated with the technology, all of which need to be considered. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**The Mark of the Assassin**
The Mark of the Assassin:
The Mark of the Assassin is a 1998 spy fiction novel by Daniel Silva.
Synopsis:
When a terrorist bomb blows Flight 002 out of the sky off the east coast, there is only one chilling clue. A body found near the crash site bears the deadly calling card of an elusive, lethal assassin - three bullets to the face. Michael Osbourne of the CIA knows the markings. Personally.
International titles:
Portuguese: A Marca do Assassino. (The Mark of the Assassin). (2011). ISBN 9789722523219 | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Postpartum depression**
Postpartum depression:
Postpartum depression (PPD), also called postnatal depression, is a type of mood disorder experienced after childbirth, which can affect both sexes. Symptoms may include extreme sadness, low energy, anxiety, crying episodes, irritability, and changes in sleeping or eating patterns. PPD can also negatively affect the newborn child.While the exact cause of PPD is unclear, the cause is believed to be a combination of physical, emotional, genetic, and social factors. These may include factors such as hormonal changes and sleep deprivation. Risk factors include prior episodes of postpartum depression, bipolar disorder, a family history of depression, psychological stress, complications of childbirth, lack of support, or a drug use disorder. Diagnosis is based on a person's symptoms. While most women experience a brief period of worry or unhappiness after delivery, postpartum depression should be suspected when symptoms are severe and last over two weeks.Among those at risk, providing psychosocial support may be protective in preventing PPD. This may include community support such as food, household chores, mother care, and companionship. Treatment for PPD may include counseling or medications. Types of counseling that have been found to be effective include interpersonal psychotherapy (IPT), cognitive behavioral therapy (CBT), and psychodynamic therapy. Tentative evidence supports the use of selective serotonin reuptake inhibitors (SSRIs).Postpartum depression affects roughly 8.9-10.1% of women in high income countries and 17.8-19.7% of women in low and middle income countries. Moreover, this mood disorder is estimated to affect 1% to 26% of new fathers. Postpartum psychosis, a more severe form of postpartum mood disorder, occurs in about 1 to 2 per 1,000 women following childbirth. Postpartum psychosis is one of the leading causes of the murder of children less than one year of age, which occurs in about 8 per 100,000 births in the United States.
Signs and symptoms:
Symptoms of PPD can occur any time in the first year postpartum. Typically, a diagnosis of postpartum depression is considered after signs and symptoms persist for at least two weeks.
Signs and symptoms:
Emotional Persistent sadness, anxiousness or "empty" mood Severe mood swings Frustration, irritability, restlessness, anger Feelings of hopelessness or helplessness Guilt, shame, worthlessness Low self-esteem Numbness, emptiness Exhaustion Inability to be comforted Trouble bonding with the baby Feeling inadequate in taking care of the baby Thoughts of self-harm or suicide Behavioral Lack of interest or pleasure in usual activities Low libido Changes in appetite Fatigue, decreased energy and motivation Poor self-care Social withdrawal Insomnia or excessive sleep Worry about harming self, baby, or partner Neurobiology fMRI studies indicate differences in brain activity between mothers with postpartum depression and those without. Mothers diagnosed with PPD tend to have less activity in the left frontal lobe and increased activity in the right frontal lobe when compared with healthy controls. They also exhibit decreased connectivity between vital brain structures, including the anterior cingulate cortex, dorsal lateral prefrontal cortex, amygdala, and hippocampus. Brain activation differences between depressed and nondepressed mothers is more pronounced when stimulated by non-infant emotional cues. Depressed mothers show greater neural activity in the right amygdala toward non-infant emotional cues as well as reduced connectivity between the amygdala and right insular cortex. Recent findings have also identified blunted activity in anterior cingulate cortex, striatum, orbitofrontal cortex, and insula in mothers with PPD when viewing images of their own infants.More robust studies on neural activation regarding PPD have been conducted with rodents than humans. These studies have allowed for greater isolation of specific brain regions, neurotransmitters, hormones, and steroids.
Signs and symptoms:
Onset and duration Postpartum depression onset usually begins between two weeks to a month after delivery. A study done at an inner-city mental health clinic has shown that 50% of postpartum depressive episodes there began prior to delivery. Therefore, in the DSM-5 postpartum depression is diagnosed under "depressive disorder with peripartum onset", in which "peripartum onset" is defined as anytime either during pregnancy or within the four weeks following delivery. PPD may last several months or even a year. Postpartum depression can also occur in women who have suffered a miscarriage. For fathers, several studies show that men experience the highest levels of postpartum depression between 3–6 months postpartum.
Signs and symptoms:
Parent-infant relationship Postpartum depression can interfere with normal maternal-infant bonding and adversely affect acute and longterm child development. Postpartum depression may lead mothers to be inconsistent with childcare. These childcare inconsistencies may include feeding routines, sleep routines, and health maintenance.In rare cases, or about 1 to 2 per 1,000, the postpartum depression appears as postpartum psychosis. In these, or among women with a history of previous psychiatric hospital admissions, infanticide may occur. In the United States, postpartum depression is one of the leading causes of annual reported infanticide incidence rate of about 8 per 100,000 births.According to research published in the American Journal of Obstetrics and Gynecology, children can experience the effects of postpartum depression. If a mother experiences postpartum depression that goes untreated, it can have adverse effects on her children. When a child is in infancy, these problems can include unusual amounts of crying (colic) and not having normal sleeping patterns. These problems can have a cyclical effect, meaning that they can further agitate the mother's postpartum depression and can even lead to the mother further developing postpartum depression. These cyclical effects can affect the way the mother maintains her relationship with her child. These can include the stopping of breastfeeding, as well as negative emotions such as withdrawal, disengagement, and even hostility. If a mother develops a hostile relationship, it can lead to extreme outcomes such as infanticide.
Signs and symptoms:
As the child grows older, postpartum depression can lead to the child experiencing irregularities in cognitive processes, behaviors, and emotions. In addition to these abnormalities, children who grew up around postpartum depression also are susceptible to developing violent tendencies.
Signs and symptoms:
Postpartum depression in fathers Paternal postpartum depression has not been studied as much as its maternal counterpart. However, postpartum depression affects 8 to 10% of fathers. In men, postpartum depression is typically defined as "an episode of major depressive disorder (MDD) occurring soon after the birth of a child". There are no set criteria for men to have postpartum depression. The cause may be distinct in males. Causes of paternal postpartum depression include hormonal changes during pregnancy, which can be indicative of father-child relationships. For instance, male depressive symptoms have been associated with low testosterone levels in men. Low prolactin, estrogen, and vasopressin levels have been associated with struggles with father-infant attachment, which can lead to depression in first-time fathers. Symptoms of postpartum depression in men are extreme sadness, fatigue, anxiety, irritability, and suicidal thoughts. Postpartum depression in men is most likely to occur 3–6 months after delivery, and is correlated with maternal depression, meaning that if the mother is experiencing postpartum depression, then the father is at a higher risk of developing the illness as well. Postpartum depression in men leads to an increase risk of suicide, while also limiting healthy infant-father attachment. Men who experience PPD can exhibit poor parenting behaviors, distress, and reduce infant interaction. Reduced paternal interaction can later lead to cognitive and behavioral problems in children. Children as young as 3.5 years old experience problems with internalizing and externalizing behaviors, indicating that paternal postpartum depression can have long-term consequences. Furthermore, if children as young as two are not frequently read to, this negative parent-child interaction can have a harmful impact on their expressive vocabulary.
Causes:
The cause of PPD is unknown. Hormonal and physical changes, personal and family history of depression, and the stress of caring for a new baby all may contribute to the development of postpartum depression.Evidence suggests that hormonal changes may play a role. Understanding the neuroendocrinology characteristic of PPD has proven to be particularly challenging given the erratic changes to the brain and biological systems during pregnancy and postpartum. A review of exploratory studies in PPD have observed that women with PPD have more dramatic changes in HPA axis activity, however directionality of specific hormone increases or decreases remain mixed. Hormones which have been studied include estrogen, progesterone, thyroid hormone, testosterone, corticotropin releasing hormone, endorphins, and cortisol. Estrogen and progesterone levels drop back to pre-pregnancy levels within 24 hours of giving birth, and that sudden change may cause it. Aberrant steroid hormone–dependent regulation of neuronal calcium influx via extracellular matrix proteins and membrane receptors involved in responding to the cell's microenvironment might be important in conferring biological risk. The use of synthetic oxytocin, a birth-inducing drug, has been linked to increased rates of postpartum depression and anxiety.Fathers, who are not undergoing profound hormonal changes, can also have postpartum depression. The cause may be distinct in males.
Causes:
Profound lifestyle changes that are brought about by caring for the infant are also frequently hypothesized to cause PPD. However, little evidence supports this hypothesis. Mothers who have had several previous children without experiencing PPD can nonetheless experience it with their latest child. Despite the biological and psychosocial changes that may accompany pregnancy and the postpartum period, most women are not diagnosed with PPD. Many mothers are unable to get the rest they need to fully recover from giving birth. Sleep deprivation can lead to physical discomfort and exhaustion, which can contribute to the symptoms of postpartum depression.
Causes:
Risk factors While the causes of PPD are not understood, a number of factors have been suggested to increase the risk. These risks can be broken down into two categories, biological and psychosocial: Biological Administration of labor-inducing medication synthetic oxytocin Chronic illnesses caused by neuroendocrine irregularities Genetic history of PPD Hormone irregularities Inflammatory illnesses (irritable bowel syndrome, fibromyalgia) Cigarette smokingThe risk factors for postpartum depression can be broken down into two categories as listed above, biological and psychosocial. Certain biological risk factors include the administration of oxytocin to induce labor. Chronic illnesses such as diabetes, or Addison's disease, as well as issues hypothalamic-pituitary-adrenal dysregulation (which controls hormonal responses), inflammatory processes like asthma or celiac disease, and genetic vulnerabilities such as a family history of depression or PPD. Chronic illnesses caused by neuroendocrine irregularities including irritable bowl syndrome and fibromyalgia typically put individuals at risk for further health complications. However, it has been found that these diseases do not increase risk for postpartum depression, these factors are known to correlate with PPD. This correlation does not mean these factors are causal. Cigarette smoking has been known to have additive effects. Some studies have found a link between PPD and low levels of DHA (an omega-3 fatty acid) in the mother. A correlation between postpartum thyroiditis and postpartum depression has been proposed but remains controversial. There may also be a link between postpartum depression and anti-thyroid antibodies.
Causes:
Psychosocial Prenatal depression or anxiety A personal or family history of depression Moderate to severe premenstrual symptoms Stressful life events experienced during pregnancy Postpartum blues Birth-related psychological trauma Birth-related physical trauma History of sexual abuse Childhood trauma Previous stillbirth or miscarriage Formula-feeding rather than breast-feeding Low self-esteem Childcare or life stress Low social support Poor marital relationship or single marital status Low socioeconomic status A lack of strong emotional support from spouse, partner, family, or friends Infant temperament problems/colic Unplanned/unwanted pregnancy Breastfeeding difficulties Maternal age, family food insecurity and violence against womenThe psychosocial risk factors for postpartum depression include severe life events, some forms of chronic strain, relationship quality, and support from partner and mother. There is a need for more research in regard to the link between psychosocial risk factors and postpartum depression. Some psychosocial risk factors can be linked to the social determinants of health. Women with fewer resources indicate a higher level of postpartum depression and stress than those women with more resources, such as financial.Rates of PPD have been shown to decrease as income increases. Women with fewer resources may be more likely to have an unintended or unwanted pregnancy, increasing risk of PPD. Women with fewer resources may also include single mothers of low income. Single mothers of low income may have more limited access to resources while transitioning into motherhood. These women already have fewer spending options, and having a child may spread those options even further. Low-income women are frequently trapped in a cycle of poverty, unable to advance, affecting their ability to access and receive quality healthcare to diagnose and treat postpartum depression.Studies have also shown a correlation between a mother's race and postpartum depression. African American mothers have been shown to have the highest risk of PPD at 25%, while Asian mothers had the lowest at 11.5%, after controlling for social factors such as age, income, education, marital status, and baby's health. The PPD rates for First Nations, Caucasian and Hispanic women fell in between.Migration away from a cultural community of support can be a factor in PPD. Traditional cultures around the world prioritize organized support during postpartum care to ensure the mother's mental and physical health, wellbeing, and recovery.One of the strongest predictors of paternal PPD is having a partner who has PPD, with fathers developing PPD 50% of the time when their female partner has PPD.Sexual orientation has also been studied as a risk factor for PPD. In a 2007 study conducted by Ross and colleagues, lesbian and bisexual mothers were tested for PPD and then compared with a heterosexual sample group. It was found that lesbian and bisexual biological mothers had significantly higher Edinburgh Postnatal Depression Scale scores than did the heterosexual women in the sample. Postpartum depression is more common among lesbian women than heterosexual women, which can be attributed to lesbian women's higher depression prevalence. Lesbian women have a higher risk of depression because they are more likely to have been treated for depression and to have attempted or contemplated suicide than heterosexual women. These higher rates of PPD in lesbian/bisexual mothers may reflect less social support, particularly from their families of origin and additional stress due to homophobic discrimination in society.There is a call to integrate both a consideration of biological and psychosocial risk factors for PPD when treating and researching the illness.
Causes:
Violence A meta-analysis reviewing research on the association of violence and postpartum depression showed that violence against women increases the incidence of postpartum depression. About one-third of women throughout the world will experience physical or sexual violence at some point in their lives. Violence against women occurs in conflict, post-conflict, and non-conflict areas. The research reviewed only looked at violence experienced by women from male perpetrators. Further, violence against women was defined as "any act of gender-based violence that results in, or is likely to result in, physical, sexual, or psychological harm or suffering to women". Psychological and cultural factors associated with increased incidence of postpartum depression include family history of depression, stressful life events during early puberty or pregnancy, anxiety or depression during pregnancy, and low social support. Violence against women is a chronic stressor, so depression may occur when someone is no longer able to respond to the violence.
Diagnosis:
Criteria Postpartum depression in the DSM-5 is known as "depressive disorder with peripartum onset". Peripartum onset is defined as starting anytime during pregnancy or within the four weeks following delivery. There is no longer a distinction made between depressive episodes that occur during pregnancy or those that occur after delivery. Nevertheless, the majority of experts continue to diagnose postpartum depression as depression with onset anytime within the first year after delivery.The criteria required for the diagnosis of postpartum depression are the same as those required to make a diagnosis of non-childbirth related major depression or minor depression. The criteria include at least five of the following nine symptoms, within a two-week period: Feelings of sadness, emptiness, or hopelessness, nearly every day, for most of the day or the observation of a depressed mood made by others Loss of interest or pleasure in activities Weight loss or decreased appetite Changes in sleep patterns Feelings of restlessness Loss of energy Feelings of worthlessness or guilt Loss of concentration or increased indecisiveness Recurrent thoughts of death, with or without plans of suicide Differential diagnosis Postpartum blues Postpartum blues, commonly known as "baby blues," is a transient postpartum mood disorder characterized by milder depressive symptoms than postpartum depression. This type of depression can occur in up to 80% of all mothers following delivery. Symptoms typically resolve within two weeks. Symptoms lasting longer than two weeks are a sign of a more serious type of depression. Women who experience "baby blues" may have a higher risk of experiencing a more serious episode of depression later on.
Diagnosis:
Psychosis Postpartum psychosis is not a formal diagnosis, but is widely used to describe a psychiatric emergency that appears to occur in about 1 in a 1000 pregnancies, in which symptoms of high mood and racing thoughts (mania), depression, severe confusion, loss of inhibition, paranoia, hallucinations and delusions begin suddenly in the first two weeks after delivery; the symptoms vary and can change quickly. It is different from postpartum depression and from maternity blues. It may be a form of bipolar disorder. It is important not to confuse psychosis with other symptoms that may occur after delivery, such as delirium. Delirium typically includes a loss of awareness or inability to pay attention.About half of women who experience postpartum psychosis have no risk factors; but a prior history of mental illness, especially bipolar disorder, a history of prior episodes of postpartum psychosis, or a family history put some at a higher risk.Postpartum psychosis often requires hospitalization, where treatment is antipsychotic medications, mood stabilizers, and in cases of strong risk for suicide, electroconvulsive therapy.The most severe symptoms last from 2 to 12 weeks, and recovery takes 6 months to a year. Women who have been hospitalized for a psychiatric condition immediately after delivery are at a much higher risk of suicide during the first year after delivery.Childbirth-Related/Postpartum Posttraumatic Stress Disorder Parents may suffer from post-traumatic stress disorder (PTSD), or suffer post-traumatic stress disorder symptoms, following childbirth. While there has been debate in the medical community as to whether childbirth should be considered a traumatic event, the current consensus is childbirth can be a traumatic event. The DSM-IV and DSM-5 (standard classifications of mental disorders used by medical professionals) do not explicitly recognize childbirth-related PTSD, but both allow childbirth to be considered as a potential cause of PTSD. Childbirth-related PTSD is closely related to postpartum depression. Research indicates mothers who have childbirth-related PTSD also commonly have postpartum depression. Childbirth-related PTSD and postpartum depression have some common symptoms. Although both diagnoses overlap in their signs and symptoms, some symptoms specific to postpartum PTSD include being easily startled, recurring nightmares and flashbacks, avoiding the baby or anything that reminds one of birth, aggression, irritability, and panic attacks. Real or perceived trauma before, during, or after childbirth is a crucial element in diagnosing childbirth-related PTSD.Currently, there are no widely recognized assessments that measure postpartum post-traumatic stress disorder in medical settings. Existing PTSD assessments (such as the DSM-IV) have been used to measure childbirth-related PTSD. Some surveys exist to measure childbirth-related PTSD specifically, however, these are not widely used outside of research settings.Approximately 3-6% of mothers in the postpartum period have childbirth-related PTSD. The percentage of individuals with childbirth-related PTSD is approximately 15-18% in high-risk samples (women who experience severe birth complications, have a history of sexual/physical violence, or have other risk factors). Research has identified several factors which increase the chance of developing childbirth-related PTSD. These include a negative subjective experience of childbirth, maternal mental health (prenatal depression, perinatal anxiety, acute postpartum depression, and history of psychological problems), history of trauma, complications with delivery and baby (for example emergency cesarean section or NICU admittance), and a low level of social support.Childbirth-related PTSD has several negative health effects. Research suggests that childbirth-related PTSD may negatively affect the emotional attachment between mother and child. However, maternal depression or other factors may also explain this negative effect. Childbirth-related PTSD in the postpartum period may also lead to issues with the child's social-emotional development. Current research suggests childbirth-related PTSD results in lower breastfeeding rates and may prevent parents from breastfeeding for the desired amount of time.
Screening:
Screening for postpartum depression is critical as up to 50% of cases go undiagnosed in the US, emphasizing the significance of comprehensive screening measures. In the US, the American College of Obstetricians and Gynecologists suggests healthcare providers consider depression screening for perinatal women. Additionally, the American Academy of Pediatrics recommends pediatricians screen mothers for PPD at 1-month, 2-month and 4-month visits. However, many providers do not consistently provide screening and appropriate follow-up. For example, in Canada, Alberta is the only province with universal PPD screening. This screening is carried out by Public Health nurses with the baby's immunization schedule. In Sweden, Child Health Services offer a free program for new parents that includes screening mothers for PPD at 2 months postpartum. However, there are concerns about adherence to screening guidelines regarding maternal mental health.The Edinburgh Postnatal Depression Scale, a standardized self-reported questionnaire, may be used to identify women who have postpartum depression. If the new mother scores 13 or more, she likely has PPD and further assessment should follow.Healthcare providers may take a blood sample to test if another disorder is contributing to depression during the screening.The Edinburgh Postnatal Depression Scale, is used within the first week of their newborn being admitted. If mothers receive a score less than 12 they are told to be reassessed because of the depression testing protocol. It is also advised that mother's in the NICU to get screened every four to six weeks as their infant remains in the neonatal intensive care unit. Mothers who score between twelve and nineteen on the EPDS are offered two types of support. The mothers are offered LV treatment provided by a nurse in the NICU and they can be referred to the mental health professional services. If a mother receives a three on item number ten of the EPDS they are immediately referred to the social work team as they may be suicidal.It is critical to acknowledge the diversity of patient populations diagnosed with postpartum depression and how this may impact the reliability of the screening tools used. There are cultural differences in how patients express symptoms of postpartum depression; those in non-western countries exhibit more physical symptoms, whereas those in western countries have more feelings of sadness. Depending on one's cultural background, symptoms of postpartum depression may manifest differently, and non-Westerners being screened in Western countries may be misdiagnosed because their screening tools do not account for cultural diversity. Aside from culture, it is also important to consider one's social context, as women with low socioeconomic status may have additional stressors that affect their postpartum depression screening scores.
Prevention:
A 2013 Cochrane review found evidence that psychosocial or psychological intervention after childbirth helped reduce the risk of postnatal depression. These interventions included home visits, telephone-based peer support, and interpersonal psychotherapy. Support is an important aspect of prevention, as depressed mothers commonly state that their feelings of depression were brought on by "lack of support" and "feeling isolated."Across different cultures, traditional rituals for postpartum care may be preventative for PPD, but are more effective when the support is welcomed by the mother.In couples, emotional closeness and global support by the partner protect against both perinatal depression and anxiety. In 2014, Alasoom and Koura found that compared to 42.9 percent of women who did not get spousal support, only 14.7 percent of women who got spousal assistance had PPD. Further factors such as communication between the couple and relationship satisfaction have a protective effect against anxiety alone.In those who are at risk counselling is recommended. In 2018, 24% of areas in the UK have no access to perinatal mental health specialist services.Preventative treatment with antidepressants may be considered for those who have had PPD previously. However, as of 2017, the evidence supporting such use is weak.
Treatments:
Treatment for mild to moderate PPD includes psychological interventions or antidepressants. Women with moderate to severe PPD would likely experience a greater benefit with a combination of psychological and medical interventions. Light aerobic exercise has been found to be useful for mild and moderate cases.
Treatments:
Therapy Both individual social and psychological interventions appear equally effective in the treatment of PPD. Social interventions include individual counseling and peer support, while psychological interventions include cognitive behavioral therapy (CBT) and interpersonal therapy (IPT). Interpersonal therapy (IPT) has shown to be effective in focusing specifically on the mother and infant bond. Support groups and group therapy options focused on psychoeducation around postpartum depression have been shown to enhance the understanding of postpartum symptoms and often assist in finding further treatment options. Other forms of therapy, such as group therapy, home visits, counseling, and ensuring greater sleep for the mother may also have a benefit. While specialists trained in providing counseling interventions often serve this population in need, results from a recent systematic review and meta-analysis found that nonspecialist providers, including lay counselors, nurses, midwives, and teachers without formal training in counseling interventions, often provide effective services related to perinatal depression and anxiety.Internet-based cognitive behavioral therapy (iCBT) has shown promising results with lower negative parenting behavior scores and lower rates of anxiety, stress, and depression. iCBT may be beneficial for mothers who have limitations in accessing in person CBT. However, the long term benefits have not been determined.
Treatments:
Medication A 2010 review found few studies of medications for treating PPD noting small sample sizes and generally weak evidence. Some evidence suggests that mothers with PPD will respond similarly to people with major depressive disorder. There is low-certainty evidence which suggests that selective serotonin reuptake inhibitors (SSRIs) are effective treatment for PPD. The first-line anti-depressant medication of choice is sertraline, an SSRI, as very little of it passes into the breast milk and, as a result, to the child. However, a recent study has found that adding sertraline to psychotherapy does not appear to confer any additional benefit. Therefore, it is not completely clear which antidepressants, if any, are most effective for treatment of PPD, and for whom antidepressants would be a better option than non-pharmacotherapy.Some studies show that hormone therapy may be effective in women with PPD, supported by the idea that the drop in estrogen and progesterone levels post-delivery contribute to depressive symptoms. However, there is some controversy with this form of treatment because estrogen should not be given to people who are at higher risk of blood clots, which include women up to 12 weeks after delivery. Additionally, none of the existing studies included women who were breastfeeding. However, there is some evidence that the use of estradiol patches might help with PPD symptoms.Oxytocin has been shown to be an effective anxiolytic and in some cases antidepressant treatment in men and women. Exogenous oxytocin has only been explored as a PPD treatment with rodents, but results are encouraging for potential application in humans.In 2019, the FDA approved brexanolone, a synthetic analog of the neurosteroid allopregnanolone, for use intravenously in postpartum depression. Allopregnanolone levels drop after giving birth, which may lead to women becoming depressed and anxious. Some trials have demonstrated an effect on PPD within 48 hours from the start of infusion. Other new allopregnanolone analogs under evaluation for use in the treatment of PPD include zuranolone and ganaxolone.Brexanolone has risks that can occur during administration, including excessive sedation and sudden loss of consciousness, and therefore has been approved under the Risk Evaluation and Mitigation Strategy (REMS) program. The mother is to enrolled prior to receiving the medication. It is only available to those at certified health care facilities with a health care provider who can continually monitor the patient. The infusion itself is a 60-hour, or 2.5 day, process. People's oxygen levels are to be monitored with a pulse oximeter. Side effects of the medication include dry mouth, sleepiness, somnolence, flushing and loss of consciousness. It is also important to monitor for early signs of suicidal thoughts or behaviors.In 2023, the FDA approved zuranolone, sold under the brand name Zurzuvae for treatment of postpartum depression. Zuranolone is administered through a pill, which is more convienient than Brexanolone, which is administered through an intravenous injection.
Treatments:
Breastfeeding Government guidance recommends caution when administering antidepressant medications during breastfeeding. Most antidepressants are excreted in breast milk in low quantities which can have adverse effect on babies. Regarding allopregnanolone, very limited data did not indicate a risk for the infant.
Other Electroconvulsive therapy (ECT) has shown efficacy in women with severe PPD that have either failed multiple trials of medication-based treatment or cannot tolerate the available antidepressants. Tentative evidence supports the use of repetitive transcranial magnetic stimulation (rTMS).As of 2013 it is unclear if acupuncture, massage, bright lights, or taking omega-3 fatty acids are useful.
Resources:
International Postpartum Support International is the most recognized international resource for those with PPD as well as healthcare providers. It brings together those experiencing PPD, volunteers, and professionals to share information, referrals, and support networks. Services offered by PSI include the website (with support, education, and local resource info), coordinators for support and local resources, online weekly video support groups in English and Spanish, free weekly phone conference with chats with experts, educational videos, closed Facebook groups for support, and professional training of healthcare workers.
Resources:
United States Educational interventions Educational interventions can help women struggling with postpartum depression (PPD) to cultivate coping strategies and develop resiliency. The phenomenon of "scientific motherhood" represents the origin of women's education on perinatal care with publications like Ms. circulating some of the first press articles on PPD that helped to normalize the symptoms that women experienced. Feminist writings on PPD from the early seventies shed light on the darker realities of motherhood and amplified the lived experiences of mothers with PPD.
Resources:
Instructional videos have been popular among women who turn to the internet for PPD treatment, especially when the videos are interactive and get patients involved in their treatment plan. Since the early 2000s, video tutorials on PPD have been integrated into many web-based training programs for individuals with PPD and are often considered a type of evidence-based management strategy for individuals. This can take the form of objective-based learning, detailed exploration of case studies, resource guides for additional support and information, etc.
Resources:
Government-funded programs The National Child and Maternal Health Education Program functions as a larger education and outreach program supported by the National Institute of Child Health and Human Development (NICHD) and the National Institute of Health. The NICHD has worked alongside organizations like the World Health Organization to conduct research on the psychosocial development of children with part of their efforts going towards the support of mothers' health and safety. Training and education services are offered through the NICHD to equip women and their health care providers with evidence-based knowledge on PPD.Other initiatives include the Substance Abuse and Mental Health Services Administration (SAMHSA) whose disaster relief program provides medical assistance at both the national and local level. The disaster relief fund not only helps to raise awareness of the benefits of having healthcare professionals screen for PPD, but also helps childhood professionals (home visitors and early care providers) develop the skills to diagnose and prevent PPD. The Infant and Early Childhood Mental Health Consultation (IECMH) center is a related technical assistance program that utilizes evidence-based treatments services in order to address issues of PPD. The IECMH facilitates parenting and home visit programs, early care site interventions with parents and children and a variety of other consultation-based services. The IECMH's initiatives seek to educate home visitors on screening protocols for PPD as well as ways to refer depressed mothers to professional help.
Resources:
Links to government-funded programs [1] [2] [3] [4] Psychotherapy Therapeutic methods of intervention can begin as early as a few days post-birth when most mothers are discharged from hospitals. Research surveys have revealed a paucity of professional, emotional support for women struggling in the weeks following delivery despite there being a heightened risk for PPD for new mothers during this transitional period.
Resources:
Community-based support A lack of social support has been identified as a barrier to seeking help for postpartum depression. Peer-support programs have been identified as an effective intervention for women experiencing symptoms of postpartum depression. In-person, online, and telephone support groups are available to both women and men throughout the United States. Peer-support models are appealing to many women because they are offered in a group and outside of the mental-health setting. The website Postpartum Progress provides a comprehensive list of support groups separated by state and includes the contact information for each group. The National Alliance on Mental Illness lists a virtual support group titled "The Shades of Blue project," which is available to all women via the submission of a name and email address. Additionally, NAMI recommends the website "National Association of Professional and Peer Lactation Supports of Color" for mothers in need of a lactation supporter. Lactation assistance is available either online or in-person, if there is support nearby.
Resources:
Personal narratives & memoirs Postpartum Progress is a blog focused on being a community of mothers talking openly about postpartum depression and other mental health conditions associated. Story-telling and online communities reduce the stigma around PPD and promote peer-based care. Postpartum Progress is specifically relevant to people of color and queer folks due to an emphasis on cultural competency.
Resources:
Hotlines & telephone interviews Hotlines, chat lines, and telephone interviews offer immediate, emergency support for those experiencing PPD. Telephone-based peer support can be effective in the prevention and treatment of postpartum depression among women at high-risk. Established examples of telephone hotlines include: National Alliance on Mental Illness: 800-950-NAMI (6264), National Suicide Prevention Lifeline: 800-273-TALK (8255), Postpartum Support International: 800-944-4PPD (4773), and SAMHSA's National Hotline: 1-800-662-HELP (4357). Postpartum Health Alliance has an immediate, 24/7 support line in San Diego/San Diego Access and Crisis Line at (888) 724–7240, in which you can talk with mothers who have recovered from PPD and trained providers.However, hotlines can lack cultural competency which is crucial in quality healthcare, specifically for people of color. Calling the police or 911, specifically for mental health crises, is dangerous for many people of color. Culturally and structurally competent emergency hotlines are a huge need in PPD care.
Resources:
National Alliance on Mental Illness: 800-950-NAMI (6264) National Suicide Prevention Lifeline: 800-273-TALK (8255) Postpartum Support International: 800-944-4PPD (4773) SAMHSA's National Hotline: 1-800-662-HELP (4357) Self-care & well-being activities Women demonstrated an interest in self-care and well-being in an online PPD prevention program. Self-care activities, specifically music therapy, are accessible to most communities and valued among women as a way to connect with their children and manage symptoms of depression. Well-being activities associated with being outdoors, including walking and running, were noted amongst women as a way to help manage mood.
Resources:
Accessibility to care Those with PPD come across many help-seeking barriers, including lack of knowledge, stigma about symptoms, as well as health service barriers. There are also attitudinal barriers to seeking treatment, including stigma. Interpersonal relationships with friends and family, as well as institutional and financial obstacles serve as help-seeking barriers. The history of mistrust within the United States healthcare system or negative health experiences can influence one's willingness and adherence to seek postpartum depression treatment. Cultural responses must be adequate in PPD healthcare and resources. Representation and cultural competency are crucial in equitable healthcare for PPD. Different ethic groups may believe that healthcare providers will not respect their cultural values or religious practices, which influences their willingness to use mental health services or be prescribed antidepressant medications. Additionally, resources for PPD are limited and often don't incorporate what mothers would prefer. The use of technology can be a beneficial way to deliver mothers with resources because it is accessible and convenient.
Epidemiology:
North America United States Within the United States, the prevalence of postpartum depression was lower than the global approximation at 11.5% but varied between states from as low as 8% to as high as 20.1%. The highest prevalence in the US is found among women who are American Indian/Alaska Natives or Asian/Pacific Islanders, possess less than 12 years of education, are unmarried, smoke during pregnancy, experience over two stressful life events, or who's full term infant is low-birthweight or was admitted to a Newborn Intensive Care Unit. While US prevalence decreased from 2004 to 2012, it did not decrease among American Indian/Alaska Native women or those with full term, low-birthweight infants.Even with the variety of studies, it is difficult to find the exact rate as approximately 60% of US women are not diagnosed and of those diagnosed approximately 50% are not treated for PPD. Cesarean section rates did not affect the rates of PPD. While there is discussion of postpartum depression in fathers, there is no formal diagnosis for postpartum depression in fathers.
Epidemiology:
Canada Canada has one of the largest refugee resettlement in the world with an equal percentage of women to men. This means that Canada has a disproportionate percentage of women that develop post-partum depression since there is an increased risk among the refugee population. In a blind study, where women had to reach out and participate, around 27% of the sample population had symptoms consistent with post-partum depression without even knowing. Also found that on average 8.46 women had minor and major PPDS was found to be 8.46 and 8.69% respectively. The main factors that were found to contribute in this study were the stress during pregnancy, the availability of support after, and a prior diagnosis of depression were all found to be factors. Canada has the specific population demographics that also involve a large amount of immigrant and indigenous women which creates a specific cultural demographic localized to Canada. In this study researchers found that these two populations were at significantly higher risk compared to "Canadian born non-indigenous mothers". This study found that risk factors such as low education, low income cut off, taking antidepressants, and low social support are all factors that contribute to the higher percentage of these population in developing PPDS. Specifically, indigenous mothers had the most risk factors then immigrant mothers with non-indigenous Canadian women being closer to the overall population.
Epidemiology:
South America A main issue surrounding PPD is the lack of study and the lack of reported prevalence that is based on studies developed in Western economically developed countries. In countries such as Brazil, Guyana, Costa Rica, Italy, Chile, and South Africa there is actually a prevalence of report, around 60%. In an itemized research analysis put a mean prevalence at 10-15% percent but explicitly stated that cultural factors such as perception of mental health and stigma could possibly be preventing accurate reporting. The analysis for South America shows that PPD occurs at a high rate looking comparatively at Brazil (42%) Chile (4.6-48%) Guyana and Colombia (57%) and Venezuela (22%). In most of these countries PPD is not considered a serious condition for women and therefore there is an absence of support programs for prevention and treatment in health systems. Specifically, in Brazil PPD is identified through the family environment whereas in Chile PPD manifests itself through suicidal ideation and emotional instability. In both cases most women feel regret and refuse to take care of the child showing that this illness is serious for both the mother and child.
Epidemiology:
Asia From a selected group of studies found from a literature search, researchers discovered many demographic factors of Asian populations that showed significant association with PPD. Some of these include age of mother at the time of childbirth as well as older age at marriage. Being a migrant and giving birth to a child overseas has also been identified as a risk factor for PPD. Specifically for Japanese women who were born and raised in Japan but who gave birth to their child in Hawaii, USA, about 50% of them experienced emotional dysfunction during their pregnancy. In fact, all women who gave birth for the first time who were included in the study experienced PPD. In immigrant Asian Indian women, the researchers found a minor depressive symptomatology rate of 28% and an additional major depressive symptomatology rate of 24% likely due to different health care attitudes in different cultures and distance form family leading to homesickness.In the context of Asian countries, premarital pregnancy is an important risk factor for PPD. This is because it is considered highly unacceptable in most Asian culture as there is a highly conservative attitude toward sex among Asian people than people in the west. In addition, conflicts between mother and daughter-in-law are notoriously common in Asian societies as traditionally for them, marriage means the daughter-in-law joining and adjusting to the groom's family completely. These conflicts may be responsible for emergence of PPD. Regarding gender of the child, many studies have suggested dissatisfaction in infant's gender (birth of a baby girl) is a risk factor for PPD. This is because in some Asian cultures, married couples are expected by the family to have at least one son to maintain the continuity of the bloodline which might lead a woman to experience PPD if she cannot give birth to a baby boy.
Epidemiology:
Europe There is a general assumption that Western cultures are homogenous and that there are no significant differences in psychiatric disorders across Europe and the USA. However, in reality factors associated with maternal depression, including work and environmental demands, access to universal maternity leave, health care, and financial security, are regulated and influenced by local policies that differ across countries. For example, European social policies differ from country-to-country contrary to the US, all countries provide some form of paid universal maternity leave and free health care. Studies also found differences in symptomatic manifestations of PPD between European and American women. Women from Europe reported higher scores of anhedonia, self-blaming, and anxiety, while women from the US disclosed more severe insomnia, depressive feelings, and thoughts of self-harming. Additionally, there are differences in prescribing patterns and attitudes towards certain medications between the US and Europe which are indicative of how different countries approach treatment, and their different stigmas.
Epidemiology:
Africa Africa, like all other parts of the world struggles with a burden of postpartum depression. Current studies estimate the prevalence to be 15-25% but this is likely higher due to a lack of data and recorded cases. The magnitude of postpartum depression in South Africa is between 31.7% and 39.6%, in Morocco between 6.9% and 14%, in Nigeria between 10.7% and 22.9%, in Uganda 43%, in Tanzania 12%, in Zimbabwe 33%, in Sudan 9.2%, in Kenya between 13% and 18.7% and, 19.9% for participants in Ethiopia according to studies carried out in these countries among postpartum mothers between the ages of 17–49. This demonstrates the gravity of this problem in Africa, and the need for postpartum depression to be taken seriously as a public health concern in the continent. Additionally, each of these studies were conducted using Western developed assessment tools. Cultural factors can affect diagnosis and can be a barrier to assessing the burden of disease. Some recommendations to combat postpartum depression in Africa include considering postpartum depression as a public health problem that is neglected among postpartum mothers. Investing in research to assess the actual prevalence of postpartum depression, and encourage early screening, diagnosis and treatment of postpartum depression as an essential aspect of maternal care throughout Africa.Issues in Reporting Prevalence Most studies regarding PPD are done using self-report screenings which are less reliable than clinical interviews. This use of self-report may have results that underreport symptoms and thus postpartum depression rates.
History:
Prior to the 19th century Western medical science's understanding and construction of postpartum depression has evolved over the centuries. Ideas surrounding women's moods and states have been around for a long time, typically recorded by men. In 460 B.C., Hippocrates wrote about puerperal fever, agitation, delirium, and mania experienced by women after child birth. Hippocrates' ideas still linger in how postpartum depression is seen today.A woman who lived in the 14th century, Margery Kempe, was a Christian mystic. She was a pilgrim known as "Madwoman" after having a tough labor and delivery. There was a long physical recovery period during which she started descending into "madness" and became suicidal. Based on her descriptions of visions of demons and conversations she wrote about that she had with religious figures like God and the Virgin Mary, historians have identified what Margery Kempe was experiencing as "postnatal psychosis" and not postpartum depression. This distinction became important to emphasize the difference between postpartum depression and postpartum psychosis. A 16th century physician, Castello Branco, documented a case of postpartum depression without the formal title as a relatively healthy woman with melancholy after childbirth, remained insane for a month, and recovered with treatment. Although this treatment was not described, experimental treatments began to be implemented for postpartum depression for the centuries that followed. Connections between female reproductive function and mental illness would continue to center around reproductive organs from this time all the way through to modern age, with a slowly evolving discussion around "female madness".
History:
19th century and after With the 19th century came a new attitude about the relationship between female mental illness and pregnancy, childbirth, or menstruation. The famous short story, "The Yellow Wallpaper", was published by Charlotte Perkins Gilman in this period. In the story, an unnamed woman journals her life when she is treated by her physician husband, John, for hysterical and depressive tendencies after the birth of their baby. Gilman wrote the story to protest societal oppression of women as the result of her own experience as a patient.Also during the 19th century, gynecologists embraced the idea that female reproductive organs, and the natural processes they were involved in, were at fault for "female insanity." Approximately 10% of asylum admissions during this time period are connected to "puerperal insanity," the named intersection between pregnancy or childbirth and female mental illness. It wasn't until the onset of the twentieth century that the attitude of the scientific community shifted once again: the consensus amongst gynecologists and other medical experts was to turn away from the idea of diseased reproductive organs and instead towards more "scientific theories" that encompassed a broadening medical perspective on mental illness.
Society and culture:
Legal recognition Recently, postpartum depression has become more widely recognized in society. In the US, the Patient Protection and Affordable Care Act included a section focusing on research into postpartum conditions including postpartum depression. Some argue that more resources in the form of policies, programs, and health objectives need to be directed to the care of those with PPD.
Society and culture:
Role of stigma When stigma occurs, a person is labelled by their illness and viewed as part of a stereotyped group. There are three main elements of stigmas, 1) problems of knowledge (ignorance or misinformation), 2) problems of attitudes (prejudice), 3) problems of behavior (discrimination). Specifically regarding PPD, it is often left untreated as women frequently report feeling ashamed about seeking help and are concerned about being labeled as a "bad mother" if they acknowledge that they are experiencing depression. Although there has been previous research interest in depression-related stigma, few studies have addressed PPD stigma. One study studied PPD stigma through examining how an education intervention would impact it. They hypothesized that an education intervention would significantly influence PPD stigma scores. Although they found some consistencies with previous mental health stigma studies, for example, that males had higher levels of personal PPD stigma than females, most of the PPD results were inconsistent with other mental health studies. For example, they hypothesized that education intervention would lower PPD stigma scores, but in reality there was no significant impact and also familiarity with PPD was not associated with one's stigma towards people with PPD. This study was a strong starting point for further PPD research, but clearly indicates more needs to be done in order to learn what the most effective anti-stigma strategies are specifically for PPD.Postpartum depression is still linked to significant stigma. This can also be difficult when trying to determine the true prevalence of postpartum depression. Participants in studies about PPD carry their beliefs, perceptions, cultural context and stigma of mental health in their cultures with them which can affect data. The stigma of mental health - with or without support from family members and health professionals - often deters women from seeking help for their PPD. When medical help is achieved, some women find the diagnosis helpful and encourage a higher profile for PPD amongst the health professional community.
Society and culture:
Cultural beliefs Postpartum depression can be influenced by sociocultural factors. There are many examples of particular cultures and societies that hold specific beliefs about PPD. Malay culture holds a belief in Hantu Meroyan; a spirit that resides in the placenta and amniotic fluid.When this spirit is unsatisfied and venting resentment, it causes the mother to experience frequent crying, loss of appetite, and trouble sleeping, known collectively as "sakit meroyan". The mother can be cured with the help of a shaman, who performs a séance to force the spirits to leave.Some cultures believe that the symptoms of postpartum depression or similar illnesses can be avoided through protective rituals in the period after birth. These may include offering structures of organized support, hygiene care, diet, rest, infant care, and breastfeeding instruction. The rituals appear to be most effective when the support is welcomed by the mother.Some Chinese women participate in a ritual that is known as "doing the month" (confinement) in which they spend the first 30 days after giving birth resting in bed, while the mother or mother-in-law takes care of domestic duties and childcare. In addition, the new mother is not allowed to bathe or shower, wash her hair, clean her teeth, leave the house, or be blown by the wind.
Society and culture:
Media Certain cases of postpartum mental health concerns received attention in the media and brought about dialogue on ways to address and understand more on postpartum mental health. Andrea Yates, a former nurse, became pregnant for the first time in 1993. After giving birth to five children in the coming years, she had severe depression and had many depressive episodes. This led to her believing that her children needed to be saved, and that by killing them, she could rescue their eternal souls. She drowned her children one by one over the course of an hour, by holding their heads under water in their family bathtub. When called into trial, she felt that she had saved her children rather than harming them and that this action would contribute to defeating Satan.This was one of the first public and notable cases of postpartum psychosis, which helped create dialogue on women's mental health after childbirth. The court found that Yates was experiencing mental illness concerns, and the trial started the conversation of mental illness in cases of murder and whether or not it would lessen the sentence or not. It also started a dialogue on women going against "maternal instinct" after childbirth and what maternal instinct was truly defined by.Yates' case brought wide media attention to the problem of filicide, or the murder of children by their parents. Throughout history, both men and women have perpetrated this act, but study of maternal filicide is more extensive. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Wc (Unix)**
Wc (Unix):
wc (short for word count) is a command in Unix, Plan 9, Inferno, and Unix-like operating systems. The program reads either standard input or a list of computer files and generates one or more of the following statistics: newline count, word count, and byte count. If a list of files is provided, both individual file and total statistics follow.
Example:
Sample execution of wc: The first column is the count of newlines, meaning that the text file foo has 40 newlines while bar has 2294 newlines- resulting in a total of 2334 newlines. The second column indicates the number of words in each text file showing that there are 149 words in foo and 16638 words in bar – giving a total of 16787 words. The last column indicates the number of characters in each text file, meaning that the file foo has 947 characters while bar has 97724 characters – 98671 characters all in all.
Example:
Newer versions of wc can differentiate between byte and character count. This difference arises with Unicode which includes multi-byte characters. The desired behaviour is selected with the -c or -m options.
Through a pipeline, it can also be used to preview the output size of a command with a potentially large output, without it printing the text into the console:
History:
wc is part of the X/Open Portability Guide since issue 2 of 1987. It was inherited into the first version of POSIX.1 and the Single Unix Specification. It appeared in Version 1 Unix.GNU wc used to be part of the GNU textutils package; it is now part of GNU coreutils. The version of wc bundled in GNU coreutils was written by Paul Rubin and David MacKenzie.A wc command is also part of ASCII's MSX-DOS2 Tools for MSX-DOS version 2.The command is available as a separate package for Microsoft Windows as part of the GnuWin32 project and the UnxUtils collection of native Win32 ports of common GNU Unix-like utilities.The wc command has also been ported to the IBM i operating system.
Usage:
wc -c <filename> prints the byte count wc -l <filename> prints the line count wc -m <filename> prints the character count wc -w <filename> prints the word count wc -L <filename> prints the length of the longest line (GNU extension) | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Plasmid preparation**
Plasmid preparation:
A plasmid preparation is a method of DNA extraction and purification for plasmid DNA, it is an important step in many molecular biology experiments and is essential for the successful use of plasmids in research and biotechnology. Many methods have been developed to purify plasmid DNA from bacteria. During the purification procedure, the plasmid DNA is often separated from contaminating proteins and genomic DNA.
Plasmid preparation:
These methods invariably involve three steps: growth of the bacterial culture, harvesting and lysis of the bacteria, and purification of the plasmid DNA. Purification of plasmids is central to molecular cloning. A purified plasmid can be used for many standard applications, such as sequencing and transfections into cells.
Growth of the bacterial culture:
Plasmids are almost always purified from liquid bacteria cultures, usually E. coli, which have been transformed and isolated. Virtually all plasmid vectors in common use encode one or more antibiotic resistance genes as a selectable marker, for example a gene encoding ampicillin or kanamycin resistance, which allows bacteria that have been successfully transformed to multiply uninhibited. Bacteria that have not taken up the plasmid vector are assumed to lack the resistance gene, and thus only colonies representing successful transformations are expected to grow.
Growth of the bacterial culture:
Bacteria are grown under favourable conditions.
Harvesting and lysis of the bacteria:
There are several methods for cell lysis, including alkaline lysis, mechanical lysis, and enzymatic lysis.
Harvesting and lysis of the bacteria:
Alkaline lysis The most common method is alkaline lysis, which involves the use of a high concentration of a basic solution, such as sodium hydroxide, to lyse the bacterial cells. When bacteria are lysed under alkaline conditions (pH 12.0–12.5) both chromosomal DNA and protein are denatured; the plasmid DNA however, remains stable. Some scientists reduce the concentration of NaOH used to 0.1M in order to reduce the occurrence of ssDNA. After the addition of acetate-containing neutralization buffer to lower the pH to around 7, the large and less supercoiled chromosomal DNA and proteins form large complexes and precipitate; but the small bacterial DNA plasmids stay in solution.
Harvesting and lysis of the bacteria:
Mechanical lysis Mechanical lysis involves the use of physical force, such as grinding or sonication, to break down bacterial cells and release the plasmid DNA. There are several different mechanical lysis methods that can be used, including French press, bead-beating, and ultrasonication.
Harvesting and lysis of the bacteria:
Enzymatic lysis Enzymatic lysis, also called Lysozyme lysis, involves the use of enzymes to digest the cell wall and release the plasmid DNA. The most commonly used enzyme for this purpose is lysozyme, which breaks down the peptidoglycan in the cell wall of Gram-positive bacteria. Lysozyme is usually added to the bacterial culture, followed by heating and/or shaking the culture to release the plasmid DNA.
Preparations by size:
Plasmid preparation can be divided into five main categories based on the scale of the preparation: minipreparation, midipreparation, maxipreparation, megapreparation, and gigapreparation. The choice of which method to use will depend on the amount of plasmid DNA required, as well as the specific application for which it will be used.Kits are available from varying manufacturers to purify plasmid DNA, which are named by size of bacterial culture and corresponding plasmid yield. In increasing order they are: miniprep, midiprep, maxiprep, megaprep, and gigaprep. The plasmid DNA yield will vary depending on the plasmid copy number, type and size, the bacterial strain, the growth conditions, and the kit.
Preparations by size:
Minipreparation Minipreparation of plasmid DNA is a rapid, small-scale isolation of plasmid DNA from bacteria. Commonly used miniprep methods include alkaline lysis and spin-column based kits. It is based on the alkaline lysis method. The extracted plasmid DNA resulting from performing a miniprep is itself often called a "miniprep".
Minipreps are used in the process of molecular cloning to analyze bacterial clones. A typical plasmid DNA yield of a miniprep is 5 to 50 µg depending on the cell strain.
Miniprep of a large number of plasmids can also be done conveniently on filter paper by lysing the cell and eluting the plasmid on to filter paper.
Midipreparation The starting E. coli culture volume is 15-25 mL of Lysogeny broth (LB) and the expected DNA yield is 100-350 µg.
Maxipreparation The starting E. coli culture volume is 100-200 mL of LB and the expected DNA yield is 500-850 µg.
Megapreparation The starting E. coli culture volume is 500 mL – 2.5 L of LB and the expected DNA yield is 1.5-2.5 mg.
Gigapreparation The starting E. coli culture volume is 2.5-5 L of LB and the expected DNA yield is 7.5–10 mg.
Purification of plasmid DNA:
It is important to consider the downstream applications of the plasmid DNA when choosing a purification method. For example, if the plasmid is to be used for transfection or electroporation, a purification method that results in high purity and low endotoxin levels is desirable. Similarly, if the plasmid is to be used for sequencing or PCR, a purification method that results in high yield and minimal contaminants is desirable. However, multiple methods of nucleic acid purification exist. All work on the principle of generating conditions where either only the nucleic acid precipitates, or only other biomolecules precipitate, allowing the nucleic acid to be separated.
Purification of plasmid DNA:
Ethanol precipitation Ethanol precipitation is a widely used method for purifying and concentrating nucleic acids, including plasmid DNA. The basic principle of this method is that nucleic acids are insoluble in ethanol or isopropanol but soluble in water. Therefore, it works by using ethanol as an antisolvent of DNA, causing it to precipitate out of solution and then it can be collected by centrifugation. The soluble fraction is discarded to remove other biomolecules.
Purification of plasmid DNA:
Spin column Spin column-based nucleic acid purification is a method of purifying DNA, RNA or plasmid from a sample using a spin column filter. The method is based on the principle of selectively binding nucleic acids to a solid matrix in the spin column, while other contaminants, such as proteins and salts, are washed away. The conditions are then changed to elute the purified nucleic acid off the column using a suitable elution buffer.
Purification of plasmid DNA:
Phenol–chloroform extraction The basic principle of the phenol-chloroform extraction is that DNA and RNA are relatively insoluble in phenol and chloroform, while other cellular components are relatively soluble in these solvents. The addition of a phenol/chloroform mixture will dissolve protein and lipid contaminants, leaving the nucleic acids in the aqueous phase. It also denatures proteins, like DNase, which is especially important if the plasmids are to be used for enzyme digestion. Otherwise, smearing may occur in enzyme restricted form of plasmid DNA.
Purification of plasmid DNA:
Beads-based extraction In beads-based extraction, addition of a mixture containing magnetic beads commonly made of iron ions binds to plasmid DNA, separating them from unwanted compounds by a magnetic rod or stand. The plasmid-bound beads are then released by removal of the magnetic field and extracted in an elution solution for down-stream experiments such as transformation or restriction digestion. This form of miniprep can also be automated, which increases the conveniency while reducing mechanical error. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Single scan dynamic molecular imaging technique**
Single scan dynamic molecular imaging technique:
Single Scan Dynamic Molecular Imaging Technique is a positron emission tomography (PET) based neuroimaging technique that allows detection of dopamine released in the brain during a cognitive or behavioral processing. The technique was developed by a psychiatry resident Rajendra Badgaiyan and his colleagues at Massachusetts General Hospital Boston.. The technique has been used to detect dopamine released during cognitive, behavioral and emotional tasks by a number of investigators. This technique has for the first time allowed scientists to detect changes in the concentration of neurotransmitters released acutely during task performance. It expanded the scope of neuroimaging studies by allowing detection of neurochemical changes associated with the brain processing. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Well travelled road effect**
Well travelled road effect:
The well travelled road effect is a cognitive bias in which travellers will estimate the time taken to traverse routes differently depending on their familiarity with the route. Frequently travelled routes are assessed as taking a shorter time than unfamiliar routes. This effect creates errors when estimating the most efficient route to an unfamiliar destination, when one candidate route includes a familiar route, whilst the other candidate route includes no familiar routes. The effect is most salient when subjects are driving, but is still detectable for pedestrians and users of public transport. The effect has been observed for centuries but was first studied scientifically in the 1980s and 1990s following from earlier "heuristics and biases" work undertaken by Daniel Kahneman and Amos Tversky.Much like the Stroop task, it is hypothesised that drivers use less cognitive effort when traversing familiar routes and therefore underestimate the time taken to traverse the familiar route. The well travelled road effect has been hypothesised as a reason that self-reported experience curve effects are overestimated. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**MYD88**
MYD88:
Myeloid differentiation primary response 88 (MYD88) is a protein that, in humans, is encoded by the MYD88 gene.
Model organisms:
Model organisms have been used in the study of MYD88 function. The gene was originally discovered and cloned by Dan Liebermann and Barbara Hoffman in mice.
Model organisms:
In that species it is a universal adapter protein as it is used by almost all TLRs (except TLR 3) to activate the transcription factor NF-κB. Mal (also known as TIRAP) is necessary to recruit Myd88 to TLR 2 and TLR 4, and MyD88 then signals through IRAK. It also interacts functionally with amyloid formation and behavior in a transgenic mouse model of Alzheimer's disease.
Model organisms:
A conditional knockout mouse line, called Myd88tm1a(EUCOMM)Wtsi was generated as part of the International Knockout Mouse Consortium program — a high-throughput mutagenesis project to generate and distribute animal models of disease to interested scientists. Male and female animals underwent a standardized phenotypic screen to determine the effects of deletion. Twenty-one tests were carried out on homozygous mutant animals, revealing one abnormality: male mutants had an increased susceptibility to bacterial infection.
Function:
The MYD88 gene provides instructions for making a protein involved in signaling within immune cells. The MyD88 protein acts as an adapter, connecting proteins that receive signals from outside the cell to the proteins that relay signals inside the cell.
Function:
In innate immunity, the MyD88 plays a pivotal role in immune cell activation through Toll-like receptors (TLRs), which belong to large group of pattern recognition receptors (PRR). In general, these receptors sense common patterns which are shared by various pathogens – Pathogen-associated molecular pattern (PAMPs), or which are produced/released during cellular damage – damage-associated molecular patterns (DAMPs).TLRs are homologous to Toll receptors, which were first described in the onthogenesis of fruit flies Drosophila, being responsible for dorso-ventral development. Hence, TLRs have been proved in all animals from insects to mammals. TLRs are located either on the cellular surface (TLR1, TLR2, TLR4, TLR5, TLR6) or within endosomes (TLR3, TLR7, TLR8, TLR9) sensing extracellular or phagocytosed pathogens, respectively. TLRs are integral membrane glycoproteins with typical semicircular-shaped extracellular parts containing leucine-rich repeats responsible for ligand binding, and Intracellular parts containing Toll-Interleukin receptor (TIR) domain.After ligand binding, all TLRs apart from TLR3, interact with adaptor protein MyD88. Another adaptor protein, which is activated by TLR3 and TLR4, is called TIR domain-containing adapter-inducing IFN-β (TRIF). Subsequently, these proteins activate two important transcription factors: NF-κB is a dimeric protein responsible for expression of various inflammatory cytokines, chemokines and adhesion and costimulatory molecules, which in turn triggers acute inflammation and stimulation of adaptive immunity IRFs is a group of proteins responsible for expression of type I interferons setting the so-called antiviral state of a cell.TLR7 and TLR9 activate both NF-κB and IRF3 through MyD88-dependent and TRIF-independent pathway, respectively.The human ortholog MYD88 seems to function similarly to mice, since the immunological phenotype of human cells deficient in MYD88 is similar to cells from MyD88 deficient mice. However, available evidence suggests that MYD88 is dispensable for human resistance to common viral infections and to all but a few pyogenic bacterial infections, demonstrating a major difference between mouse and human immune responses. Mutation in MYD88 at position 265 leading to a change from leucine to proline have been identified in many human lymphomas including ABC subtype of diffuse large B-cell lymphoma and Waldenström's macroglobulinemia.
Interactions:
Myd88 has been shown to interact with:
Gene polymorphisms:
Various single nucleotide polymorphisms (SNPs) of the MyD88 have been identified. and for some of them an association with susceptibility to various infectious diseases and to some autoimmune diseases like ulcerative colitis was found. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Cartan's theorems A and B**
Cartan's theorems A and B:
In mathematics, Cartan's theorems A and B are two results proved by Henri Cartan around 1951, concerning a coherent sheaf F on a Stein manifold X. They are significant both as applied to several complex variables, and in the general development of sheaf cohomology.
Cartan's theorems A and B:
Theorem B is stated in cohomological terms (a formulation that Cartan (1953, p. 51) attributes to J.-P. Serre): Analogous properties were established by Serre (1957) for coherent sheaves in algebraic geometry, when X is an affine scheme. The analogue of Theorem B in this context is as follows (Hartshorne 1977, Theorem III.3.7): These theorems have many important applications. For instance, they imply that a holomorphic function on a closed complex submanifold, Z, of a Stein manifold X can be extended to a holomorphic function on all of X. At a deeper level, these theorems were used by Jean-Pierre Serre to prove the GAGA theorem.
Cartan's theorems A and B:
Theorem B is sharp in the sense that if H1(X, F) = 0 for all coherent sheaves F on a complex manifold X (resp. quasi-coherent sheaves F on a noetherian scheme X), then X is Stein (resp. affine); see (Serre 1956) (resp. (Serre 1957) and (Hartshorne 1977, Theorem III.3.7)). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Carbonates on Mars**
Carbonates on Mars:
Head (vessel) Evidence for carbonates on Mars was first discovered in 2008. Previously, most remote sensing instruments such as OMEGA and THEMIS—sensitive to infrared emissivity spectral features of carbonates—had not suggested the presence of carbonate outcrops, at least at the 100 m or coarser spatial scales available from the returned data.Though ubiquitous, a 2003 study of carbonates on Mars showed that they are dominated by Magnesite (MgCO3) in Martian dust, had mass fractions less than 5%, and could have formed under current atmospheric conditions. Furthermore, with the exception of the surface dust component, by 2007 carbonates had not been detected by any in situ mission, even though mineralogic modeling did not preclude small amounts of calcium carbonate in Independence class rocks of Husband Hill in Gusev crater (note: An IAU naming convention within Gusev is not yet established).
Remote sensing data:
The first successful identification of a strong infrared spectral signature from surficial carbonate minerals of local scale (< 10 km²) was made by the MRO-CRISM team. Spectral modeling in 2007 identified a key deposit in Nili Fossae dominated by a single mineral phase that was spatially associated with olivine outcrops. The dominant mineral appeared to be magnesite, while morphology inferred with HiRISE and thermal properties suggested that the deposit was lithic. Stratigraphically, this layer appeared between phyllosilicates below and mafic cap rocks above, temporally between the Noachian and Hesperian eras. Even though infrared spectra are representative of minerals to less than ≈0.1 mm depths (in contrast to gamma spectra which are sensitive to tens of cm depths), stratigraphic, morphologic, and thermal properties are consistent with the existence of the carbonate as outcrop rather than alteration rinds. Nevertheless, the morphology was distinct from typical terrestrial sedimentary carbonate layers suggesting formation from local aqueous alteration of olivine and other igneous minerals. However, key implications were that the alteration would have occurred under moderate pH and that the resulting carbonates were not exposed to sustained low pH aqueous conditions even as recently as the Hesperian. This increased the likelihood of local and regional scale geologic conditions on Mars that were favorable to analogs of terrestrial biological activity over geologically significant intervals.As of 2012, the absence of more extensive carbonate deposits on Mars was thought by some scientists to be due to global dominance of low pH aqueous environments.
Remote sensing data:
Even the least soluble carbonate, siderite (FeCO3), precipitates only at a pH greater than 5.Evidence for significant quantities of carbonate deposits on the surface began to increase in 2008 when the Thermal and Evolved Gas Analyzer (TEGA) and WCL experiments on the 2007 Phoenix Mars lander found between 3–5wt% calcite (CaCO3) and an alkaline soil. In 2010 analyses by the Mars Exploration Rover Spirit, identified outcrops rich in magnesium-iron carbonate (16–34 wt%) in the Columbia Hills of Gusev crater, most likely precipitated from carbonate-bearing solutions under hydrothermal conditions at near-neutral pH in association with volcanic activity during the Noachian era.After Spirit Rover stopped working scientists studied old data from the Miniature Thermal Emission Spectrometer, or Mini-TES and confirmed the presence of large amounts of carbonate-rich rocks, which means that regions of the planet may have once harbored water. The carbonates were discovered in an outcrop of rocks called "Comanche."Carbonates (calcium or iron carbonates) were discovered in a crater on the rim of Huygens Crater, located in the Iapygia quadrangle. The impact on the rim exposed material that had been dug up from the impact that created Huygens. These minerals represent evidence that Mars once had a thicker carbon dioxide atmosphere with abundant moisture. These kind of carbonates only form when there is a lot of water. They were found with the Compact Reconnaissance Imaging Spectrometer for Mars (CRISM) instrument on the Mars Reconnaissance Orbiter. Earlier, the instrument had detected clay minerals. The carbonates were found near the clay minerals. Both of these minerals form in wet environments. It is supposed that billions of years age Mars was much warmer and wetter. At that time, carbonates would have formed from water and the carbon dioxide-rich atmosphere. Later the deposits of carbonate would have been buried. The double impact has now exposed the minerals. Earth has vast carbonate deposits in the form of limestone. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Bicycle collecting**
Bicycle collecting:
As with many consumer products, early bicycles were purchased solely for their usefulness or fashionableness and discarded as they wore out or were replaced by newer models. Some items were thrown into storage and survived, but many others went to the scrapyard. Decades later, those with an interest in cycling and history began to seek out older bikes, collecting different varieties. Like other forms of collecting, bike collectors can be completists or specialists, and many have extensive holdings in bike parts or literature, in addition to complete bicycles.
North America:
Due to the tremendous number of bicycle manufacturers and models that have appeared over the past 150 years, most collectors specialize in a particular style or period of bicycles. Currently, there are three primary periods of particular collector interest in North America, although many collectors will further specialize in the products of a single manufacturer or even examples of a single model within a given period. The major periods are: High Wheel and Antique (Early 19th century-1933)—Early bicycles were all experiments and came in a dizzying variety of shapes. From primitive “hobby horses” to the giant High Wheel or Penny Farthing bicycles of the 1880s, collectors have gathered and studied these strange designs. Although many of these models are extremely rare, their peculiar shapes are fascinating and offer insight into the development of mechanical solutions that eventually resulted in the fairly standardized “safety bicycle” of the 1890s. Also included in this category are the early safety bicycles, which featured wooden rim wheels, skinny tires, and slightly larger wheel diameter than what became standard later.
North America:
Balloon Tire Classics (1933–1965)—This period is dominated by the cruiser style bicycles of Schwinn and other manufacturers. These bikes featured wide balloon tires and heavy frames, for improved durability. The children’s market was a focus during this era, leading to elaborate streamline styling and loads of accessories: lights, speedometers, springer (suspension) forks, horns, luggage racks, and more. These bikes were neglected and abused until the mid 1970s when Leon Dixon began penning a series of articles for magazines such as Popular Mechanics and organized the earliest collector swap meets. Soon prices of old cruisers began to rise. Today, this is probably the most popular area of bike collecting.
North America:
Wheelie bikes and Early BMX (1965–1980)—This fast-growing segment of the hobby in North America focuses on the Schwinn Sting-Rays, Raleigh Choppers and other banana seat bikes of the 1960s and the early BMX models that grew out of them. The Sting-Rays offered a huge assortment of accessories, much like the old cruisers, but over this period the bikes were stripped down and made stronger and stronger to withstand the rigors of dirt track racing and trick riding. Prices in this category have begun to rise recently as the children of the 1960s reach the age where they have the money, the time, and the inclination to collect. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**MIR1269A**
MIR1269A:
MicroRNA 1269a is a protein that in humans is encoded by the MIR1269A gene.
Function:
microRNAs (miRNAs) are short (20-24 nt) non-coding RNAs that are involved in post-transcriptional regulation of gene expression in multicellular organisms by affecting both the stability and translation of mRNAs. miRNAs are transcribed by RNA polymerase II as part of capped and polyadenylated primary transcripts (pri-miRNAs) that can be either protein-coding or non-coding. The primary transcript is cleaved by the Drosha ribonuclease III enzyme to produce an approximately 70-nt stem-loop precursor miRNA (pre-miRNA), which is further cleaved by the cytoplasmic Dicer ribonuclease to generate the mature miRNA and antisense miRNA star (miRNA*) products. The mature miRNA is incorporated into a RNA-induced silencing complex (RISC), which recognizes target mRNAs through imperfect base pairing with the miRNA and most commonly results in translational inhibition or destabilization of the target mRNA. The RefSeq represents the predicted microRNA stem-loop. [provided by RefSeq, Sep 2009]. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**DOTA (chelator)**
DOTA (chelator):
DOTA (also known as tetraxetan) is an organic compound with the formula (CH2CH2NCH2CO2H)4. The molecule consists of a central 12-membered tetraaza (i.e., containing four nitrogen atoms) ring. DOTA is used as a complexing agent, especially for lanthanide ions. Its complexes have medical applications as contrast agents and cancer treatments.
Terminology:
The acronym DOTA (for dodecane tetraacetic acid) is shorthand for both the tetracarboxylic acid and its various conjugate bases. In the area of coordination chemistry, the tetraacid is called H4DOTA and its fully deprotonated derivative is DOTA4−. Many related ligands are referred to using the DOTA acronym, although these derivatives are generally not tetracarboxylic acids or the conjugate bases.
Structure:
DOTA is derived from the macrocycle known as cyclen. The four secondary amine groups are modified by replacement of the N-H centers with N-CH2CO2H groups. The resulting aminopolycarboxylic acid, upon ionization of the carboxylic acid groups, is a high affinity chelating agent for di- and trivalent cations. The tetracarboxylic acid was first reported in 1976. At the time of its discovery DOTA exhibited the largest known formation constant for the complexation (chelating) of Ca2+ and Gd3+ ions. Modified versions of DOTA were first reported in 1988 and this area has proliferated since.As a polydentate ligand, DOTA envelops metal cations, but the denticity of the ligand depends on the geometric tendencies of the metal cation. The main applications involve the lanthanides and in such complexes DOTA functions as an octadentate ligand, binding the metal through four amine and four carboxylate groups. Most such complexes feature an additional water ligand, giving an overall coordination number of nine.For most transition metals, DOTA functions as a hexadentate ligand, binding through the four nitrogen and two carboxylate centres. The complexes have octahedral coordination geometry, with two pendent carboxylate groups. In the case of [Fe(DOTA)]−, the ligand is heptadentate.
Uses:
Cancer treatment and diagnosis DOTA can be conjugated to monoclonal antibodies by attachment of one of the four carboxyl groups as an amide. The remaining three carboxylate anions are available for binding to the yttrium ion. The modified antibody accumulates in the tumour cells, concentrating the effects of the radioactivity of 90Y. Drugs containing this module receive an International Nonproprietary Name ending in tetraxetan: Yttrium (90Y) clivatuzumab tetraxetan Yttrium (90Y) tacatuzumab tetraxetanDOTA can also be linked to molecules that have affinity for various structures. The resulting compounds are used with a number of radioisotopes in cancer therapy and diagnosis (for example in positron emission tomography).
Uses:
Affinity for somatostatin receptors, which are found on neuroendocrine tumours:DOTATOC, DOTA-(Tyr3)-octreotide or edotreotide DOTA-TATE or DOTA-(Tyr3)-octreotate Affinity for the proteins streptavidin and avidin, which can be targeted at tumours by aid of monoclonal antibodies:DOTA-biotin Contrast agent The complex of Gd3+ and DOTA is used as a gadolinium-based MRI contrast agent under the name gadoteric acid.
Synthesis:
DOTA was first synthesized in 1976 from cyclen and bromoacetic acid. This method is simple and still in use. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Protonophore**
Protonophore:
A protonophore, also known as a proton translocator, is an ionophore that moves protons across lipid bilayers or other type of membranes. This would otherwise not occur as protons cations (H+) have positive charge and hydrophilic properties, making them unable to cross without a channel or transporter in the form of a protonophore. Protonophores are generally aromatic compounds with a negative charge, that are both hydrophobic and capable of distributing the negative charge over a number of atoms by π-orbitals which delocalize a proton's charge when it attaches to the molecule. Both the neutral and the charged protonophore can diffuse across the lipid bilayer by passive diffusion and simultaneously facilitate proton transport. Protonophores uncouple oxidative phosphorylation via a decrease in the membrane potential of the inner membrane of mitochondria. They stimulate mitochondria respiration and heat production. Protonophores (uncouplers) are often used in biochemistry research to help explore the bioenergetics of chemiosmotic and other membrane transport processes. It has been reported that the protonophore has antibacterial activity by perturbing bacterial proton motive force.Representative anionic protonophores include: 2,4-dinitrophenol Carbonyl cyanide-p-trifluoromethoxyphenylhydrazone (FCCP) Carbonyl cyanide m-chlorophenyl hydrazone (CCCP)Representative cationic protonophores include: C4R1 (a short-chain alkyl derivative of rhodamine 19) EllipticineRepresentative zwitterionic protonophores include: mitoFluo (10-[2-(3-hydroxy-6-oxo-xanthen-9-yl)benzoyl]oxydecyl-triphenyl-phosphonium bromide) PP6 (2-(2-Hydroxyaryl)hexylphosphonium bromide)
Mechanism of action:
The facilitated transport of protons across the biological membrane by anionic protonophore is achieved as follows.
The anionic form of the protonophore (P−) is adsorbed onto one side (Positive) of the biological membrane.
Protons (H+) from the aqueous solution combine with the anion (P−) to produce the neutral form (PH) PH diffuses across the biological membrane and dissociates into H+ and P− on the other side.
This H+ is released from the biological membrane into the other aqueous solution P− returns to the first side of the biological membrane by electrophoresis (its electrostatic attraction to the positive side of the membrane). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Suspensory behavior**
Suspensory behavior:
Suspensory behaviour is a form of arboreal locomotion or a feeding behavior that involves hanging or suspension of the body below or among tree branches. This behavior enables faster travel while reducing path lengths to cover more ground when travelling, searching for food and avoiding predators. Different types of suspensory behaviour include brachiation, climbing, and bridging. These mechanisms allow larger species to distribute their weight among smaller branches rather than balancing above these weak supports. Primates and sloths are most commonly seen using these behaviours, however, other animals such as bats may be seen hanging below surfaces to obtain food or when resting.
Biomechanics:
In primates and sloths Animals who exhibit suspensory behaviour have similar mechanisms to perform this action and often involve many different parts of their body like the trunk, shoulders and many other features of their upper body. Typically, these animals have an overall dorso-ventral flattening, a shortened lumbar region and a mediolateral expansion of the rib cage causing the scapula to be repositioned dorsally and humeral articulation to be oriented more cranially than the usual lateral placement shown in quadrupedal animals. The scapula is also longer, giving these animals a particular arm and shoulder shape. Combined, these morphologies allow for the infraspinatus muscle to be repositioned creating more resistance to trans articular tensile stress for suspending below a branch. These animals also have longer clavicles, creating a bigger projection of the shoulder which increases the ability to move when the forearm is raised above the head. To help with supporting their weight, the forelimbs are elongated. The humerus is longer as well and this helps with the movement of the deltoid muscles in the shoulder joint when the arm is moving away from the body. The triceps branchii is small and there is a shorter distance to the elbow joint and a shorter olecranon process which allows for a greater elbow extension.
Biomechanics:
There are many different ways animals, especially primates, position themselves during suspensory behaviour and these positions require different bones and muscles. Below is a list of different positions and their mechanisms.
Sit/forelimb-suspend: Most of the weight of the body is put on the ischium, however the abducted forelimbs grab a hold of a branch overhead and allow for the body to be stabilized and supports some of this weight that is being put on the ischia.
Cling/forelimb-suspend: One of the forelimbs is hanging causing more than half the weight to be put the hindlimbs and the clinging forelimb.
Forelimb-suspend: One or both arms is holding all the weight of the animal as it hangs from a branch.Unimanual forelimb-suspend: Suspension using one hand with lack of support from the rest of the body. The humerus is abducted and the elbow is usually extended completely.
Bimanual forelimb-suspend: suspensions using both hands.
Biomechanics:
Forelimb-suspend/sit: This is similar to sit/forelimb-suspend except more than half the weight is held by the forelimbs and not the ishchia. The arms of the animal are extended and the remainder of the weight is supported by the ishchia and/or feet. In this position one arm can hang, creating most of the weight to be held by the single forelimb.
Biomechanics:
Forelimb-suspend/squat: suspension from above but the lower body is in a squat position.
Forelimb-suspension/stand: Half of the weight is supported by the two forelimbs that are extended, the other half is supported from standing.
Forelimb-suspend/cling: Hind limbs are flexed and grasping a support while one or both of the forelimbs are grasping the support as well, distributing the weight evenly.
Forelimb-suspend/lie: suspension of the forelimbs with the back in a horizontal position, as if they were lying on their back.
Trunk-vertical-suspend: One or both forelimbs and one or both hind limbs carry the weight. The foot/feet are above the level of the hip. This differs from other forms as all four limbs have tension.
Unimanual flexed-elbow-suspend: Suspension with the humerus adducted and the elbow not extended. These parts of the body hold the animals entire weight.
Bimanual flexed-elbow-suspend: similar to unimanual flexed-elbow-suspension, expect both hands are involved, not just one.
Forelimb-hindlimb-suspend: hanging from the arm and foot.Ipsilateral forelimb-hind-limb-suspend: suspension with a forelimb and hind limb on the same side of the body.
Contralateral forelimb-hind-limb-suspend: suspension with a forelimb and a hind limb on the opposite sides of the body.
Tail-suspend: suspension from the tail, with no support from the rest of the body.Tail/forelimb-suspend: Half of the weight is on the tail and the other half on the forelimb(s).
Tail/hind limb-suspend: Half of the weight is on the tail and the other half on the hind limb(s).
Pronograde tail/quadrumanous-suspend: All five limbs help support the body while the back is horizontal.
Orthograde tail/quadrumanous-suspend: All five limbs help support the body while the back is vertical.
Hind limb-suspend: Suspension from the foot/feet, no support from any other body parts.
Flexed-hind limb-suspend: Knee and the hip are flexed during suspension.
Extended-hind limb-suspend: Knee and hip are extended during suspension.
Biomechanics:
In bats Roosting is a vertical upside down behaviour seen in bats which involves the use of the feet to grasp a surface. The hind limbs are very important as they provide most of the strength to support the bat. The forelimbs can be used as well, having all four limbs supporting the animal. The head and neck are usually kept at a 90° or 180° angle.
Locomotion:
Suspensory locomotion aids with reducing path lengths and covering longer distances by moving faster through branches and trees above. The movements of involved in suspensory behavior can be described as being seen most often among monkeys. The swinging motion of grabbing branch after branch with alternating hands or launching the body from one support to another losing contact with the support is very common and the most popular form of locomotion among suspensory animals. Some animals such as the platyrrhines, use their tails for traveling and usually never use their forelimbs for transportation, while some species use both their tails and forelimbs. Suspensory behavior is advantageous for avoiding predators. The quick motions and ability to escape high above the ground enables an avoidance strategy, maintaining survival. While this type of locomotion can be beneficial there can be some consequences when dealing with extreme heights as vigorously moving through the trees allows for more opportunity for injury. The easiest way for animals to avoid this consequence is using their abilities to focus on uninterrupted travel, accuracy and avoiding alternative routes.
Locomotion:
Types of locomotion Brachiation Brachiation involves the animal swinging from branch to branch in a sequence motion above the ground in a canopy of trees. Typically these movements involve both arms without the aid of the legs or tail. Tail and hind limb suspension can be used in different situations like feeding or escaping predators during drastic situations, however the use of the arms is preferred for this type of movement.
Locomotion:
Climbing Climbing consists of moving up or down a vertical surface using all four arms and legs to help move the body upward or downward. There are many different ways in which in animal can climb such as using alternating arms and legs, climbing sideways, fire-pole slides and head or bottom first decline. Vertical climbing is the most costly form of locomotion as the animal must defy gravity and move up the tree. This is particularly harder for animals with a larger body mass, as carrying their entire weight becomes more difficult with size. Also involved with climbing is a "pulling up" motion in which the animal will pull itself above a branch using both of its arms and the hind limbs launch over the branch using a swinging motion.
Locomotion:
Bridging Animals use this type of behavior when crossing between trees and other surfaces. This movement requires the use of the hind limbs to leap across extended areas. Small animals have an easier time leaping between gaps, while larger animals are more cautious due to their weight and typically swing from branch to branch instead.
Feeding:
Suspensory behaviour is very important for animals in regards to feeding. It has been reported that suspensory movements make up approximately 25% of all feeding strategies shown in primates. Suspension helps them reach fruits and other vegetation that might be difficult to obtain on foot, while allowing them to cover a large distance at a greater speed. Often in arboreal regions, flowers, fruits and other plants are located on small terminal branches and suspension enables animals to access this food while saving time and energy. By suspending below the branch they avoid a greater chance at the branch breaking and are able to keep a steady balance. Hanging by the tail is very common when foraging which permits the use of the hands and arms to not only grab food but to catch themselves if they were to slip or fall. Suspension allows for fast travel, which is helpful when collecting food as well. Speed allows animals to minimize competition while avoiding predators to ensure they grab as much food as they can in a short period of time. If an animal is in a high tree, they often eat their food then and there to avoid injury and predators. Quadrupedalism and bipedalism combined with suspensory mechanisms are crucial for providing support during feeding so the animal does not fall and risk losing the food, or risking its life. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Downcasting**
Downcasting:
In class-based programming, downcasting or type refinement is the act of casting a reference of a base class to one of its derived classes.
In many programming languages, it is possible to check through type introspection to determine whether the type of the referenced object is indeed the one being cast to or a derived type of it, and thus issue an error if it is not the case.
In other words, when a variable of the base class (parent class) has a value of the derived class (child class), downcasting is possible.
Some languages, such as OCaml, disallow downcasting.
Uses:
Downcasting is useful when the type of the value referenced by the Parent variable is known and often is used when passing a value as a parameter. In the below example, the method objectToString takes an Object parameter which is assumed to be of type String.
In this approach, downcasting prevents the compiler from detecting a possible error and instead causes a run-time error.
Uses:
Downcasting myObject to String ('(String)myObject') was not possible at compile time because there are times that myObject is String type, so only at run time can we figure out whether the parameter passed in is logical. While we could also convert myObject to a compile-time String using the universal java.lang.Object.toString(), this would risk calling the default implementation of toString() where it was unhelpful or insecure, and exception handling could not prevent this.
Uses:
In C++, run-time type checking is implemented through dynamic_cast. Compile-time downcasting is implemented by static_cast, but this operation performs no type check. If it is used improperly, it could produce undefined behavior.
Considerations:
A popular example of a badly considered design is containers of top types, like the Java containers before Java generics were introduced, which requires downcasting of the contained objects so that they can be used again. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Guaijaverin**
Guaijaverin:
Guaijaverin is the 3-O-arabinoside of quercetin. It is found in the leaves of Psidium guajava, the common guava. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Goin' Nuts**
Goin' Nuts:
Goin' Nuts is a pinball machine that was designed by Adolf Seitz, Jr. for Gottlieb in 1983. The game never went into production and only 10 prototypes were built.
Description:
The game is unique in that it has no outlanes. Also, the game starts as 3-ball multiball with auto-plunger instead of a plunger; once the player started the game all three balls are released and objective is to knock down the set of drop targets in order to score points and build up time. Rather than counting balls, a player's score is determined based on the amount of time left on the timer which counts down when there's only 1 ball left.Disadvantages of this pinball machine include damage of the playfield and toys by multiple balls nicking each other. A good player may also build up too much time which leads to lower income for the machine owner.
Design team:
Game Design: Adolf Seitz Jr.
Artwork: unidentified artist at Advertising Posters of Chicago
Digital versions:
Goin' Nuts is available as a licensed table of The Pinball Arcade for several platforms and the game was also included in the Pinball Hall of Fame: The Gottlieb Collection for the Nintendo Wii system. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Diisopromine**
Diisopromine:
Diisopromine or disoprominum, usually as the hydrochloride salt, is a synthetic spasmolytic which neutralizes spastic conditions of the biliary tract and of the sphincter of Oddi. It was discovered at Janssen Pharmaceutica in 1955. It is sold in South Africa under the brand name Agofell syrup as a mixture with sorbitol, and elsewhere as Megabyl. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Townhouse**
Townhouse:
A townhouse, townhome, town house, or town home, is a type of terraced housing. A modern townhouse is often one with a small footprint on multiple floors. In a different British usage, the term originally referred to any type of city residence (normally in London) of someone whose main or largest residence was a country house.
History:
Historically, a townhouse was the city residence of a noble or wealthy family, who would own one or more country houses in which they lived for much of the year. From the 18th century, landowners and their servants would move to a townhouse during the social season (when major balls took place).
Europe:
United Kingdom In the United Kingdom, most townhouses are terraced. Only a small minority of them, generally the largest, were detached, but even aristocrats whose country houses had grounds of hundreds or thousands of acres often lived in terraced houses in town. For example, the Duke of Norfolk owned Arundel Castle in the country, while his London house, Norfolk House, was a terraced house in St James's Square over 100 feet (30 meters) wide.
North America:
U.S. and Canada In the United States and Canada, a townhouse has two connotations. The older predates the automobile and denotes a house on a small footprint in a city, but because of its multiple floors (sometimes six or more), it has a large living space, often with servants' quarters. The small footprint of the townhouse allows it to be within walking or mass-transit distance of business and industrial areas of the city, yet luxurious enough for wealthy residents of the city.Townhouses are expensive where detached single-family houses are uncommon, such as in New York City, Chicago, Boston, Philadelphia, Montreal, Washington, D.C., and San Francisco.
North America:
Rowhouses are similar and consist of several adjacent, uniform units originally found in older, pre-automobile urban areas such as Baltimore, Philadelphia, Richmond, Virginia, Charleston, South Carolina, Savannah, Georgia and New Orleans, but now found in lower-cost housing developments in suburbs as well. A rowhouse is where there is a continuous roof and foundation, and a single wall divides adjacent townhouses, but some have a double wall with inches-wide air space in between on a common foundation. A rowhouse will generally be smaller and less luxurious than a dwelling called a townhouse.
North America:
The name townhouse or townhome was later used to describe non-uniform units in suburban areas that are designed to mimic detached or semi-detached homes. Today, the term townhouse is used to describe units mimicking a detached home that are attached in a multi-unit complex. The distinction between living units called apartments and those called townhouses is that townhouses usually consist of multiple floors and have their own outside door as opposed to having only one level and/or having access via an interior corridor hallway or via an exterior balcony-style walkway (more common in the warmer climates). Another distinction is that in most areas of the US outside of the very largest cities, apartment refers to rental housing, and townhouse typically refers to an individually owned dwelling, with no other unit beneath or above although the term townhouse-style (rental) apartment is also heard for bi-level apartments.
North America:
Townhouses can also be "stacked". Such homes have multiple units vertically (typically two), normally each with its own private entrance from the street or at least from the outside. They can be side by side in a row of three or more, in which case they are sometimes referred to as rowhouses. A townhouse in a group of two could be referred to as a townhouse, but in Canada and the US, it is typically called a semi-detached home and in some areas of western Canada, a half-duplex.
North America:
In Canada, single-family dwellings, be they any type, such as single-family detached homes, apartments, mobile homes, or townhouses, for example, are split into two categories of ownership: Condominium (strata title), where one owns the interior of the unit and also a specified share of the undivided interest of the remainder of the building and land known as common elements.
North America:
Freehold, where one owns exclusively the land and the dwelling without any condominium aspects. These may share the foundation as well but have narrow air spaces between and still referred to as a townhouse.Condominium townhouses, just like condominium apartments, are often referred to as condos, thus referring to the type of ownership rather than to the type of dwelling. Since apartment style condos are the most common, when someone refers to a condo, many erroneously assume that it must be an apartment-style dwelling and that only apartment-style dwellings can be condos. All types of dwellings can be condos, and this is therefore true of townhouses. A brownstone townhouse is a particular variety found in New York.
Asia, Australia, South Africa, Zimbabwe:
In Asia, Australia, South Africa and Zimbabwe, the usage of the term follows the North American sense. Townhouses are generally found in complexes. Large complexes often have high security, resort facilities such as swimming pools, gyms, parks and playground equipment. Typically, a townhouse has a strata title; i.e., a type of title where the common property (landscaped area, public corridors, building structure, etc.) is owned by a corporation of individual owners and the houses on the property are owned by the individual owners.
Asia, Australia, South Africa, Zimbabwe:
In population-dense Asian cities dominated by high-rise residential apartment blocks, such as Hong Kong, townhouses in private housing developments remain almost exclusively populated by the very wealthy due to the rarity and relatively large sizes of the units. Prominent examples in Hong Kong include Severn 8, in which a 5,067-square-foot (470.7 m2) townhouse sold for HK$285 million (US$37 million) in 2008, or HK$57,000 (US$7,400) per square foot, a record in Asia, and The Beverly Hills, which consists of multiple rows of townhouses with some units as large as 11,000 square feet (1,000 m2). Commonly in the suburbs of major cities, an old house on a large block of land is demolished and replaced by a short row of townhouses, built 'end on' to the street for added privacy. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Laminar organization**
Laminar organization:
A laminar organization describes the way certain tissues, such as bone membrane, skin, or brain tissues, are arranged in layers.
Types:
Embryo The earliest forms of laminar organization are shown in the diploblastic and triploblastic formation of the germ layers in the embryo. In the first week of human embryogenesis two layers of cells have formed, an external epiblast layer (the primitive ectoderm), and an internal hypoblast layer (primitive endoderm). This gives the early bilaminar disc. In the third week in the stage of gastrulation epiblast cells invaginate to form endoderm, and a third layer of cells known as mesoderm. Cells that remain in the epiblast become ectoderm. This is the trilaminar disc and the epiblast cells have given rise to the three germ layers.
Types:
Brain In the brain a laminar organization is evident in the arrangement of the three meninges, the membranes that cover the brain and spinal cord. These membranes are the dura mater, arachnoid mater, and pia mater. The dura mater has two layers a periosteal layer near to the bone of the skull, and a meningeal layer next to the other meninges.The cerebral cortex, the outer neural sheet covering the cerebral hemispheres can be described by its laminar organization, due to the arrangement of cortical neurons into six distinct layers.
Types:
Eye The eye in mammals has an extensive laminar organization. There are three main layers – the outer fibrous tunic, the middle uvea, and the inner retina. These layers have sublayers with the retina having ten ranging from the outer choroid to the inner vitreous humor and including the retinal nerve fiber layer.
Skin The human skin has a dense laminar organization. The outer epidermis has four or five layers. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Scaffolding**
Scaffolding:
Scaffolding, also called scaffold or staging, is a temporary structure used to support a work crew and materials to aid in the construction, maintenance and repair of buildings, bridges and all other man-made structures. Scaffolds are widely used on site to get access to heights and areas that would be otherwise hard to get to. Unsafe scaffolding has the potential to result in death or serious injury. Scaffolding is also used in adapted forms for formwork and shoring, grandstand seating, concert stages, access/viewing towers, exhibition stands, ski ramps, half pipes and art projects.
Scaffolding:
There are six main types of scaffolding used worldwide today. These are tube and coupler (fitting) components, prefabricated modular system scaffold components, H-frame / façade modular system scaffolds, suspended scaffolds, timber scaffolds and bamboo scaffolds (particularly in China and India). Each type is made from several components which often include: A base jack or plate which is a load-bearing base for the scaffold.
Scaffolding:
The standard, the upright component with connector joins.
The ledger, a horizontal brace.
The transom, a horizontal cross-section load-bearing component which holds the batten, board, or decking unit.
Brace diagonal and/or cross section bracing component.
Batten or board decking component used to make the working platform.
Coupler, a fitting used to join components together.
Scaffold tie, used to tie in the scaffold to structures.
Scaffolding:
Brackets, used to extend the width of working platforms.Specialized components used to aid in their use as a temporary structure often include heavy duty load bearing transoms, ladders or stairway units for the ingress and egress of the scaffold, beams ladder/unit types used to span obstacles and rubbish chutes used to remove unwanted materials from the scaffold or construction project.
History:
Antiquity Sockets in the walls around the paleolithic cave paintings at Lascaux, suggest that a scaffold system was used for painting the ceiling, over 17,000 years ago.
The Berlin Foundry Cup depicts scaffolding in ancient Greece (early 5th century BC). Egyptians, Nubians and Chinese are also recorded as having used scaffolding-like structures to build tall buildings. Early scaffolding was made of wood and secured with rope knots.
History:
Modern era Scaffolding was erected by individual firms with wildly varying standards and sizes. The process was revolutionized by Daniel Palmer Jones and David Henry Jones. Modern day scaffolding standards, practices and processes can be attributed to these men and their companies: Rapid Scaffold Tie Company Ltd, Tubular Scaffolding Company and Scaffolding Great Britain Ltd (SGB).David Palmer-Jones patented the "Scaffixer", a coupling device far more robust than rope which revolutionized scaffolding construction. In 1913, his company was commissioned for the reconstruction of Buckingham Palace, during which his Scaffixer gained much publicity. Palmer-Jones followed this up with the improved "Universal Coupler" in 1919 - this soon became the industry standard coupling and has remained so to this day.Advancements in metallurgy throughout the early 20th century saw the introduction of tubular steel water pipes (instead of timber poles) with standardized dimensions, allowing for the industrial interchangeability of parts and improving the structural stability of the scaffold. The use of diagonal bracings also helped to improve stability, especially on tall buildings. The first frame system was brought to market by SGB in 1944 and was used extensively for the postwar reconstruction.
Today:
The European Standard, BS EN 12811-1, specifies performance requirements and methods of structural and general design for access and working scaffolds. Requirements given are for scaffold structures that rely on the adjacent structures for stability. In general these requirements also apply to other types of working scaffolds.
Today:
The purpose of a working scaffold is to provide a safe working platform and access suitable for work crews to carry out their work. The European Standard sets out performance requirements for working scaffolds. These are substantially independent of the materials of which the scaffold is made. The standard is intended to be used as the basis for enquiry and design.
Today:
Materials The basic components of scaffolding are tubes, couplers and boards.
Today:
The basic lightweight tube scaffolding that became the standard and revolutionised scaffolding, becoming the baseline for decades, was invented and marketed in the mid-1950s. With one basic 24 pound unit a scaffold of various sizes and heights could be assembled easily by a couple of labourers without the nuts or bolts previously needed.Tubes are usually made either of steel or aluminium. Composite scaffolding uses filament-wound tubes of glass fibre in a nylon or polyester matrix. Because of the high cost of composite tube, it is usually only used when there is a risk from overhead electric cables that cannot be isolated. Steel tubes are either 'black' or galvanised. The tubes come in a variety of lengths and a standard outside diameter of 48.3 mm. (1.5 NPS pipe). The chief difference between the two types of metal tubes is the lower weight of aluminium tubes (1.7 kg/m as opposed to 4.4 kg/m). Aluminium tube is more flexible and has a lower resistance to stress. Tubes are generally bought in 6.3 m lengths and can then be cut down to certain typical sizes. Most large companies will brand their tubes with their name and address in order to deter theft.
Today:
Boards provide a working surface for scaffold users. They are seasoned wood and come in three thicknesses (38 mm (usual), 50 mm and 63 mm) are a standard width (225 mm) and are a maximum of 3.9 m long. The board ends are protected either by metal plates called hoop irons or sometimes nail plates, which often have the company name stamped into them. Timber scaffold boards in the UK should comply with the requirements of BS 2482. As well as timber, steel or aluminium decking is used, as well as laminate boards. In addition to the boards for the working platform, there are sole boards which are placed beneath the scaffolding if the surface is soft or otherwise suspect, although ordinary boards can also be used. Another solution, called a scaffpad, is made from a rubber base with a base plate moulded inside; these are desirable for use on uneven ground since they adapt, whereas sole boards may split and have to be replaced.
Today:
Couplers are the fittings which hold the tubes together. The most common are called scaffold couplers, and there are three basic types: right-angle couplers, putlog couplers and swivel couplers. To join tubes end-to-end joint pins (also called spigots) or sleeve couplers are used. Only right angle couplers and swivel couplers can be used to fix tube in a 'load-bearing connection'. Single couplers are not load-bearing couplers and have no design capacity.
Today:
Other common scaffolding components include base plates, ladders, ropes, anchor ties, reveal ties, gin wheels, sheeting, etc.
Most companies will adopt a specific colour to paint the scaffolding with, in order that quick visual identification can be made in case of theft. All components that are made from metal can be painted but items that are wooden should never be painted as this could hide defects.
Despite the metric measurements given, many scaffolders measure tubes and boards in imperial units, with tubes from 21 feet down and boards from 13 ft down.
Bamboo scaffolding is widely used in Hong Kong and Macau, with nylon straps tied into knots as couplers. In India, bamboo or other wooden scaffolding is also mostly used, with poles being lashed together using ropes made from coconut hair (coir).
Today:
Basic scaffolding The key elements of the scaffolding are the standard, ledger and transoms. The standards, also called uprights, are the vertical tubes that transfer the entire weight of the structure to the ground where they rest on a square base plate to spread the load. The base plate has a shank in its centre to hold the tube and is sometimes pinned to a sole board. Ledgers are horizontal tubes which connect between the standards. Transoms rest upon the ledgers at right angles. Main transoms are placed next to the standards, they hold the standards in place and provide support for boards; intermediate transoms are those placed between the main transoms to provide extra support for boards. In Canada this style is referred to as "English". "American" has the transoms attached to the standards and is used less but has certain advantages in some situations.
Today:
As well as the tubes at right angles there are cross braces to increase rigidity, these are placed diagonally from ledger to ledger, next to the standards to which they are fitted. If the braces are fitted to the ledgers they are called ledger braces. To limit sway a facade brace is fitted to the face of the scaffold every 30 metres or so at an angle of 35°-55° running right from the base to the top of the scaffold and fixed at every level.
Today:
Of the couplers previously mentioned, right-angle couplers join ledgers or transoms to standards, putlog or single couplers join board bearing transoms to ledgers - Non-board bearing transoms should be fixed using a right-angle coupler. Swivel couplers are to connect tubes at any other angle. The actual joints are staggered to avoid occurring at the same level in neighbouring standards.
The spacings of the basic elements in the scaffold are fairly standard. For a general purpose scaffold the maximum bay length is 2.1 m, for heavier work the bay size is reduced to 2 or even 1.8 m while for inspection a bay width of up to 2.7 m is allowed.
Today:
The scaffolding width is determined by the width of the boards, the minimum width allowed is 600 mm but a more typical four-board scaffold would be 870 mm wide from standard to standard. More heavy-duty scaffolding can require 5, 6 or even up to 8 boards width. Often an inside board is added to reduce the gap between the inner standard and the structure.
Today:
The lift height, the spacing between ledgers, is 2 m, although the base lift can be up to 2.7 m. The diagram above also shows a kicker lift, which is just 150 mm or so above the ground.
Today:
Transom spacing is determined by the thickness of the boards supported, 38 mm boards require a transom spacing of no more than 1.2 m while a 50 mm board can stand a transom spacing of 2.6 m and 63 mm boards can have a maximum span of 3.25 m. The minimum overhang for all boards is 50 mm and the maximum overhang is no more than 4x the thickness of the board.
Today:
Foundations Good foundations are essential. Often scaffold frameworks will require more than simple base plates to safely carry and spread the load. Scaffolding can be used without base plates on concrete or similar hard surfaces, although base plates are always recommended. For surfaces like pavements or tarmac base plates are necessary. For softer or more doubtful surfaces sole boards must be used, beneath a single standard a sole board should be at least 1,000 square centimetres (160 in2) with no dimension less than 220 millimetres (8.7 in), the thickness must be at least 35 millimetres (1.4 in). For heavier duty scaffold much more substantial baulks set in concrete can be required. On uneven ground steps must be cut for the base plates, a minimum step size of around 450 millimetres (18 in) is recommended. A working platform requires certain other elements to be safe. They must be close-boarded, have double guard rails and toe and stop boards. Safe and secure access must also be provided.
Today:
Ties Scaffolds are only rarely independent structures. To provide stability for a scaffolding (at left) framework ties are generally fixed to the adjacent building/fabric/steelwork.
Today:
General practice is to attach a tie every 4 m on alternate lifts (traditional scaffolding). Prefabricated System scaffolds require structural connections at all frames - i.e. 2–3 m centres (tie patterns must be provided by the System manufacturer/supplier). The ties are coupled to the scaffold as close to the junction of standard and ledger (node point) as possible. Due to recent regulation changes, scaffolding ties must support +/- loads (tie/butt loads) and lateral (shear) loads.
Today:
Due to the different nature of structures there is a variety of different ties to take advantage of the opportunities.
Through ties are put through structure openings such as windows. A vertical inside tube crossing the opening is attached to the scaffold by a transom and a crossing horizontal tube on the outside called a bridle tube. The gaps between the tubes and the structure surfaces are packed or wedged with timber sections to ensure a solid fit.
Today:
Box ties are used to attach the scaffold to suitable pillars or comparable features. Two additional transoms are put across from the lift on each side of the feature and are joined on both sides with shorter tubes called tie tubes. When a complete box tie is impossible a l-shaped lip tie can be used to hook the scaffold to the structure, to limit inward movement an additional transom, a butt transom, is placed hard against the outside face of the structure.
Today:
Sometimes it is possible to use anchor ties (also called bolt ties), these are ties fitted into holes drilled in the structure. A common type is a ring bolt with an expanding wedge which is then tied to a node point.
Today:
The least 'invasive' tie is a reveal tie. These use an opening in the structure but use a tube wedged horizontally in the opening. The reveal tube is usually held in place by a reveal screw pin (an adjustable threaded bar) and protective packing at either end. A transom tie tube links the reveal tube to the scaffold. Reveal ties are not well regarded, they rely solely on friction and need regular checking so it is not recommended that more than half of all ties be reveal ties.
Today:
If it is not possible to use a safe number of ties rakers can be used. These are single tubes attached to a ledger extending out from the scaffold at an angle of less than 75° and securely founded. A transom at the base then completes a triangle back to the base of the main scaffold.
Bamboo scaffolding:
Bamboo scaffolding is a type of scaffolding made from bamboo and widely used in construction work for centuries. Many famous landmarks, notably The Great Wall of China, were built using bamboo scaffolding, and its use continues today in some parts of the world.
Bamboo scaffolding:
History Bamboo scaffolding was first introduced into the building industry in Hong Kong immediately after colonization in the 1800s. It was widely used in the building of houses and multi-story buildings (up to four stories high) prior to the development of metal scaffolding. It was also useful for short-term construction projects, such as framework for temporary sheds for Cantonese Opera performances.
Bamboo scaffolding:
There are three types of scaffolding in Hong Kong: Double-row scaffold; Extended Bamboo scaffolding; Shop signs of Bamboo Scaffolding.
Bamboo scaffolding:
Gradual decline In 2013, there were 1,751 registered bamboo scaffolders and roughly 200 scaffolding companies in Hong Kong. The use of bamboo scaffolding is diminishing due to shortages in labor and material. Despite the lack of labor force and material, recently safety issues have become another serious concern.The labor shortage may be due to the reluctance of younger generations to become scaffolders. “They even think that it’s a dirty and dangerous job. They are not going to do that kind of work,” said Yu Hang Flord, who has been a scaffolder for 30 years and later became the director of Wui Fai Holdings, a member of the Hong Kong and Kowloon Scaffolders General Merchants Association. “They refuse to step in, although we give them high pay. They are scared of it. Young generations do not like jobs that involve hard work.” Another reason fewer people are becoming scaffolders is that new recruits need to undergo training with the Hong Kong Construction Industry Council in order to acquire a license. Older scaffolders generally learned in apprenticeships, and may have been able to gather more hands-on experience.Material shortages are also a contributing factor to the decline. The bamboo scaffolding material was imported from mainland China. Bamboo—which matures after three years to the wide diameter and thick skin perfect for scaffolding—came from the Shaoxing area in Guangdong. Over the past two decades, firms have had to look to Guangxi instead. The industry's fear is that one day supplies will be blocked due to export embargoes and environmental concerns. Attempts to import bamboo from Thailand, or switch to synthetic or plastic bamboo, have so far proved unsuccessful.
Bamboo scaffolding:
In many African countries, notably Nigeria, bamboo scaffolding is still used for small scale construction in urban areas. In rural areas, the use of bamboo scaffolding for construction is common. In fact, bamboo is an essential building and construction commodity in Nigeria; the bamboo materials are transported on heavy trucks and trailers from rural areas (especially the tropical rain forest) to cities and the northern part of Nigeria.
Bamboo scaffolding:
Some of the structures in relaxation and recreation centres, both in urban and rural areas of Nigeria, are put in place using bamboo materials. This is not for reasons of poverty (especially in the cities) but to add more aesthetics to these centres. Bamboo materials are still used in the construction of some bukas (local restaurants) in rural areas.
Specifications Forms of bamboo scaffolding include: Double-row ScaffoldOnly double-row bamboo scaffold is allowed to be used for working at height.
Nylon MeshThe perimeter of bamboo scaffold should be covered by nylon mesh against falling objects. The lapping of nylon mesh should be at least 100 mm wide.
Access and EgressSuitable means of access should be provided from the building or ground level to the scaffold such as gangway, stairs and ladder etc.
Catch FanSloping catch fans shall be erected at a level close to the first floor and at no more than 15 metres, vertical intervals should give a minimum horizontal protection coverage of 1500 mm. Large catch fans should be erected at specific locations to protect the public and/or workers underneath.
Platform of Catch Fan or ReceptacleA suitable receptacle, covered with galvanized zinc sheet, should be provided within each catch-fan to trap falling objects.
Steel BracketSteel brackets shall be provided for supporting the standard of scaffold at about six floor intervals. The horizontal distance between steel brackets is about 3 metres.
PutlogsMild steel bars or similar materials are required to tie any structure to maintain the bamboo scaffold in its position on every floor. The distance of adjacent putlogs is about 3 to 4 metres.
Working PlatformEvery working platform must be at least 400 mm wide and closely boarded by planks. The edges of working platforms should be protected by no less than 2 horizontal bamboo members of the scaffold, at intervals between 750 mm to 900 mm and suitable toe-boards no less than 200 mm high.
Special ScaffoldAll scaffolds with a height excess of 15 metres shall be designed by an Engineer. Competent ExaminerThey should complete a formal training in bamboo scaffolding work or hold a trade test certificate on bamboo scaffolding and have at least 10 years of relevant experience.
Trained WorkerThey should complete formal training in bamboo scaffolding work or hold a trade test certificate on bamboo scaffolding and have at least 3 years of relevant experience.
Bamboo scaffolding:
Uses in construction Bamboo scaffolding is a temporary structure to support people and materials when constructing or repairing building exteriors and interiors. In bamboo scaffolding, plastic fibre straps and bamboo shoots are bound together to form a solid and secure scaffold structure without screws. Bamboo scaffolding does not need to have a foundation on the ground, as long as the scaffolding has a fulcrum for structural support.Bamboo scaffolding is mostly seen in developing Asian countries such as India, Bangladesh, Sri Lanka, and Indonesia.
Bamboo scaffolding:
Cultural use Chinese opera theatres Chinese Opera is one of the world's "Intangible Cultural Heritages". One of bamboo scaffolding's main alternative uses is in drama theatres. The flexibility and convenience of this type of scaffolding suits stages set up for temporary use and also separates the audience from the performers.
Respecting and promoting the traditional cultures of Chinese Opera, a huge event called the West Kowloon Bamboo Theatre has been held at the West Kowloon Waterfront Promenade annually since 2012.
Yu Lan Ghost Festival Stages are built from bamboo scaffolding for the live Chinese operas and Chiu Chow–style dramas performed during every Yu Lan Ghost Festival to worship ghostly ancestors.
Bamboo scaffolding:
Cheung Chau Bun Festival The bamboo tower used in the famous Bun Scrambling Competition during the Cheung Chau Bun Festival on the island of Cheung Chau is constructed out of bamboo scaffolding. Nine thousand buns, representing fortune and blessing, are supported on the fourteen-meter tall bamboo tower in front of the Pak Tai Temple. For the Piu Sik Parade, bamboo stands and racks are used to hold the young costumed performers above the crowds.
Specialty scaffolding:
Types of scaffolding covered by the Occupational Health and Safety Administration in the United States include the following categories: Pole; tube and coupler; fabricated frame (tubular welded frame scaffolds); plasterers’, decorators’, and large area scaffolds; bricklayers' (pipe); horse; form scaffolds and carpenters’ bracket scaffolds; roof brackets; outrigger; pump jacks; ladder jacks; window jacks; crawlingboards (chicken ladders); step, platform, and trestle ladder scaffolds; single-point adjustable suspension; two-point adjustable suspension (swing stages); multipoint adjustable suspension; stonesetters’ multipoint adjustable suspension scaffolds, and masons’ multipoint adjustable suspension scaffolds; catenary; float (ship); interior hung; needle beam; multilevel suspended; mobile; repair bracket scaffolds; and stilts.
Specialty scaffolding:
Gallery of scaffold types Putlog scaffold In addition to the putlog couplers (discussed above), there are also putlog tubes. These have a flattened end or have been fitted with a blade. This feature allows the end of the tube to be inserted into or rest upon the brickwork of the structure.
A putlog scaffold may also be called a bricklayer's scaffold. As such, the scaffold consists only of a single row of standards with a single ledger. The putlogs are transoms - attached to the ledger at one end but integrated into the bricks at the other.
Spacing is the same on a putlog scaffold as on a general purpose scaffold, and ties are still required.
In recent years a number of new innovations have meant an increased scope of use for scaffolding, such as ladderbeams for spanning spaces that cannot accommodate standards and the increased use of sheeting and structure to create temporary roofs.
Specialty scaffolding:
Putlog tubes can also be used vertically when drove under downward pressure into the ground, most typically in greens and fields, where approx 1/4 of the putlog tube remains exposed above ground. The purpose for this alternative method is to create a good anchoring point for additional vertical scaffolding to clamp on to, most commonly used in live events and festivals with scaffolding poles up to 21 feet high where festoon lighting, cabling and bunting can be hung from safely.
Specialty scaffolding:
Pump-jack A pump-jack is a type of portable scaffolding system. The scaffold rests on supports attached to two or more vertical posts. The user raises the scaffolding by pumping the foot pedals on the supports, like an automobile jack.
Baker staging Baker staging is a metal scaffold which is easy to assemble. Rolling platforms typically 740 millimetres (29 in) wide by 1.8 metres (6 ft) long and 1.8 metres (6 ft) tall sections which can be stacked up to three high with the use of added outriggers. The work platform height is adjustable.
X-Deck ladder scaffolding Low level scaffolding that is height adjustable. It is a hybrid ladder scaffold work platform.
Standards:
The widespread use of scaffolding systems, along with the profound importance that they earned in modern applications such as civil engineering projects and temporary structures, led to the definition of a series of standards covering a vast number of specific issues involving scaffolding. Among the standards there are: DIN 4420, a DIN standard divided in 5 parts which covers the design and detail of scaffolds, ladder scaffolds, safety requirements and standard types, materials, components, dimensions and loadbearing capacity.
Standards:
DIN 4421, a DIN standard which covers the analysis, design and construction of falsework 29 CFR Part 1926: Safety Standards for Scaffolds Used in the Construction Industry from the U.S. Occupational Safety and Health Administration (OSHA), with an accompanying "construction eTool" | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Watchdog timer**
Watchdog timer:
A watchdog timer (or simply a watchdog), sometimes called a computer operating properly timer (COP timer), is an electronic or software timer that is used to detect and recover from computer malfunctions. Watchdog timers are widely used in computers to facilitate automatic correction of temporary hardware faults, and to prevent errant or malevolent software from disrupting system operation.
Watchdog timer:
During normal operation, the computer regularly restarts the watchdog timer to prevent it from elapsing, or "timing out". If, due to a hardware fault or program error, the computer fails to restart the watchdog, the timer will elapse and generate a timeout signal. The timeout signal is used to initiate corrective actions. The corrective actions typically include placing the computer and associated hardware in a safe state and invoking a computer reboot.
Watchdog timer:
Microcontrollers often include an integrated, on-chip watchdog. In other computers the watchdog may reside in a nearby chip that connects directly to the CPU, or it may be located on an external expansion card in the computer's chassis.
Applications:
Watchdog timers are commonly found in embedded systems and other computer-controlled equipment where humans cannot easily access the equipment or would be unable to react to faults in a timely manner. In such systems, the computer cannot depend on a human to invoke a reboot if it hangs; it must be self-reliant. For example, remote embedded systems such as space probes are not physically accessible to human operators; these could become permanently disabled if they were unable to autonomously recover from faults. In robots and other automated machines, a fault in the control computer could cause equipment damage or injuries before a human could react, even if the computer is easily accessed. A watchdog timer is usually employed in cases like these.
Applications:
Watchdog timers are also used to monitor and limit software execution time on a normally functioning computer. For example, a watchdog timer may be used when running untrusted code in a sandbox, to limit the CPU time available to the code and thus prevent some types of denial-of-service attacks. In real-time operating systems, a watchdog timer may be used to monitor a time-critical task to ensure it completes within its maximum allotted time and, if it fails to do so, to terminate the task and report the failure.
Architecture and operation:
Restarting The act of restarting a watchdog timer is commonly referred to as kicking the watchdog. Kicking is typically done by writing to a watchdog control port or by setting a particular bit in a register. Alternatively, some tightly coupled watchdog timers are kicked by executing a special machine language instruction. An example of this is the CLRWDT (clear watchdog timer) instruction found in the instruction set of some PIC microcontrollers.
Architecture and operation:
In computers that are running operating systems, watchdog restarts are usually invoked through a device driver. For example, in the Linux operating system, a user space program will kick the watchdog by interacting with the watchdog device driver, typically by writing a zero character to /dev/watchdog or by calling a KEEPALIVE ioctl. The device driver, which serves to abstract the watchdog hardware from user space programs, may also be used to configure the time-out period and start and stop the timer.
Architecture and operation:
Some watchdog timers will only allow kicks during a specific time window. The window timing is usually relative to the previous kick or, if the watchdog has not yet been kicked, to the moment the watchdog was enabled. The window begins after a delay following the previous kick, and ends after a further delay. If the computer attempts to kick the watchdog before or after the window, the watchdog will not be restarted, and in some implementations this will be treated as a fault and trigger corrective action.
Architecture and operation:
Enabling A watchdog timer is said to be enabled when operating and disabled when idle. Upon power-up, a watchdog may be unconditionally enabled or it may be initially disabled and require an external signal to enable it. In the latter case, the enabling signal may be automatically generated by hardware or it may be generated under software control.
Architecture and operation:
When automatically generated, the enabling signal is typically derived from the computer reset signal. In some systems the reset signal is directly used to enable the watchdog. In others, the reset signal is delayed so that the watchdog will become enabled at some later time following the reset. This delay allows time for the computer to boot before the watchdog is enabled. Without this delay, the watchdog would timeout and invoke a subsequent reset before the computer can run its application software — the software which kicks the watchdog — and the system would become stuck in an endless cycle of incomplete reboots.
Architecture and operation:
Single-stage watchdog Watchdog timers come in many configurations, and many allow their configurations to be altered. For example, the watchdog and CPU may share a common clock signal as shown in the block diagram below, or they may have independent clock signals. A basic watchdog timer has a single timer stage which, upon timeout, typically will reset the CPU: Multistage watchdog Two or more timers are sometimes cascaded to form a multistage watchdog timer, where each timer is referred to as a timer stage, or simply a stage. For example, the block diagram below shows a three-stage watchdog. In a multistage watchdog, only the first stage is kicked by the processor. Upon first stage timeout, a corrective action is initiated and the next stage in the cascade is started. As each subsequent stage times out, it triggers a corrective action and starts the next stage. Upon final stage timeout, a corrective action is initiated, but no other stage is started because the end of the cascade has been reached. Typically, single-stage watchdog timers are used to simply restart the computer, whereas multistage watchdog timers will sequentially trigger a series of corrective actions, with the final stage triggering a computer restart.
Architecture and operation:
Time intervals Watchdog timers may have either fixed or programmable time intervals. Some watchdog timers allow the time interval to be programmed by selecting from among a few selectable, discrete values. In others, the interval can be programmed to arbitrary values. Typically, watchdog time intervals range from ten milliseconds to a minute or more. In a multistage watchdog, each timer may have its own, unique time interval.
Corrective actions:
A watchdog timer may initiate any of several types of corrective action, including maskable interrupt, non-maskable interrupt, hardware reset, fail-safe state activation, power cycling, or combinations of these. Depending on its architecture, the type of corrective action or actions that a watchdog can trigger may be fixed or programmable. Some computers (e.g., PC compatibles) require a pulsed signal to invoke a hardware reset. In such cases, the watchdog typically triggers a hardware reset by activating an internal or external pulse generator, which in turn creates the required reset pulses.In embedded systems and control systems, watchdog timers are often used to activate fail-safe circuitry. When activated, the fail-safe circuitry forces all control outputs to safe states (e.g., turns off motors, heaters, and high-voltages) to prevent injuries and equipment damage while the fault persists. In a two-stage watchdog, the first timer is often used to activate fail-safe outputs and start the second timer stage; the second stage will reset the computer if the fault cannot be corrected before the timer elapses.
Corrective actions:
Watchdog timers are sometimes used to trigger the recording of system state information—which may be useful during fault recovery—or debug information (which may be useful for determining the cause of the fault) onto a persistent medium. In such cases, a second timer—which is started when the first timer elapses—is typically used to reset the computer later, after allowing sufficient time for data recording to complete. This allows time for the information to be saved, but ensures that the computer will be reset even if the recording process fails.
Corrective actions:
For example, the above diagram shows a likely configuration for a two-stage watchdog timer. During normal operation the computer regularly kicks Stage1 to prevent a timeout. If the computer fails to kick Stage1 (e.g., due to a hardware fault or programming error), Stage1 will eventually timeout. This event will start the Stage2 timer and, simultaneously, notify the computer (by means of a non-maskable interrupt) that a reset is imminent. Until Stage2 times out, the computer may attempt to record state information, debug information, or both. As a last resort, the computer will be reset upon Stage2 timeout.
Fault detection:
A watchdog timer provides automatic detection of catastrophic malfunctions that prevent the computer from kicking it. However, computers often have other, less-severe types of faults which do not interfere with kicking, but which still require watchdog oversight. To support these, a computer system is typically designed so that its watchdog timer will be kicked only if the computer deems the system functional. The computer determines whether the system is functional by conducting one or more fault detection tests and will kick the watchdog only if all tests have passed.In computers that are running an operating system and multiple processes, a single, simple test might be insufficient to guarantee normal operation, as it could fail to detect a subtle fault condition and therefore allow the watchdog to be kicked even though a fault condition exists. For example, in the case of the Linux operating system, a user-space watchdog daemon may simply kick the watchdog periodically without performing any tests. As long as the daemon runs normally, the system will be protected against serious system crashes such as a kernel panic. To detect less severe faults, the daemon can be configured to perform tests that cover resource availability (e.g., sufficient memory and file handles, reasonable CPU time), evidence of expected process activity (e.g., system daemons running, specific files being present or updated), overheating, and network activity, and system-specific test scripts or programs can also be run.Upon discovery of a failed test, the computer may attempt to perform a sequence of corrective actions under software control, culminating with a software-initiated reboot. If the software fails to invoke a reboot, the watchdog timer will timeout and invoke a hardware reset. In effect, this is a multistage watchdog timer in which the software comprises the first and intermediate timer stages and the hardware reset the final stage. In a Linux system, for example, the watchdog daemon could attempt to perform a software-initiated restart, which can be preferable to a hardware reset as the file systems will be safely unmounted and fault information will be logged. It is essential, however, to have the insurance provided by a hardware timer, since a software restart can fail under a number of fault conditions. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Calibre (software)**
Calibre (software):
Calibre (, stylised calibre) is a cross-platform free and open-source suite of e-book software. Calibre supports organizing existing e-books into virtual libraries, displaying, editing, creating and converting e-books, as well as syncing e-books with a variety of e-readers. Editing books is supported for EPUB and AZW3 formats. Books in other formats like MOBI must first be converted to those formats, if they are to be edited.
History:
On 31 October 2006, when Sony introduced its PRS-500 e-reader, Kovid Goyal started developing libprs500, aiming mainly to enable use of the PRS-500 formats on Linux. With support from the MobileRead forums, Goyal reverse-engineered the proprietary Broad Band eBook (BBeB) file format. In 2008, the program, for which a graphical user interface was developed, was renamed "calibre", displayed in all lowercase.
Features:
Calibre supports many file formats and reading devices. Most e-book formats can be edited, for example, by changing the font, font size, margins, and metadata, and by adding an auto-generated table of contents. Conversion and editing are easily applied to appropriately licensed digital books, but commercially purchased e-books may need to have digital rights management (DRM) restrictions removed. Calibre does not natively support DRM removal, but may allow DRM removal after installing plug-ins with such a function.Calibre allows users to sort and group e-books by metadata fields. Metadata can be pulled from many different sources, e.g., ISBNdb.com; online booksellers; and providers of free e-books and periodicals in the US and elsewhere, such as the Internet Archive, Munsey's Magazine, and Project Gutenberg; and social networking sites for readers, such as Goodreads and LibraryThing. It is possible to search the Calibre library by various fields, such as author, title, or keyword. Full text search is available from Calibre 6.0 onwards.E-books can be imported into the Calibre library, either by sideloading files manually or by wirelessly syncing an e-book reading device with the cloud storage service in which the Calibre library is backed up, or with the computer on which Calibre resides. Also, online content can be harvested and converted to e-books. This conversion is facilitated by so-called recipes, short programs written in a Python-based domain-specific language. E-books can then be exported to all supported reading devices via USB, Calibre's integrated mail server, or wirelessly. Mailing e-books enables, for example, sending personal documents to the Amazon Kindle family of e-readers and tablet computers.This can be accomplished via a web browser, if the host computer is running and the device and host computer share the same network; in this case, pushing harvested content from content sources is supported on a regular interval (called 'subscription'). Also, if the Calibre library on the host computer is stored in a cloud service, such as Box.net, Google Drive, or Dropbox, then either the cloud service or a third-party app, such as Calibre Cloud or CalibreBox, can remotely access the library.Since version 1.15, released in December 2013, Calibre also contains an application to create and edit e-books directly, similar to the more full-featured editor tools of the Sigil application, but without the latter's WYSIWYG editing mode.
Associated apps:
Calibre Cloud (free) and Calibre Cloud Pro (paid), apps by Intrepid Logic that let one "access your Calibre e-book library from anywhere in the world. Place your calibre library in your Dropbox, Box, or Google Drive folder, and be able to view, search, and download from your library anywhere". As Jane Litte at Dear Author and John Jeremy at Teleread observe: This tool can be used to "create [one's] own Cloud of eBooks" and thereby read and allow downloads and emails from one's Calibre library via the Calibre folder in Box.net, Dropbox, or Google Drive. Because the Calibre-generated local wireless feed (OPDS) can only be accessed on devices sharing the same network as the Calibre library, this feature of the Calibre Cloud apps is particularly useful when away from one's home network, because it allows one to download and read the contents of one's Calibre library via the Calibre folder in Box, Dropbox, or Google Drive.
Associated apps:
Calibre Companion (paid), an app by MultiPie, Ltd., but now abandonware. Was recommended by calibre's developers, "brings complete integration with calibre on your desktop, giving you total control over book management on your device." John Jermey at Teleread notes this app can manage Calibre/device libraries as if one's mobile device were plugged into computer; however, unlike Calibre Cloud, Calibre Companion requires users to be at a computer and use the Calibre-generated local wireless feed (OPDS).
Associated apps:
Calibre Library (paid), an app by Tony Maro that allows one to "Connect wirelessly to your Calibre e-book library or other Stanza source. Browse and download your e-books on the go." This app's operations and benefits are similar to those offered by Calibre Cloud.
Associated apps:
Calibre Sync (free), an app by Seng Jea Lee that "seamlessly connects to your Calibre Library and shows up as a connected device on Calibre. If Auto-Connect option is enabled, your device will attempt to connect to the Calibre Library when it is within the home Wi-Fi network. This allows Calibre to automatically update your device with the latest newspaper or magazines you have scheduled for download!" As with Calibre Companion, this app requires the device to be on the same network as the Calibre library.
Associated apps:
CalibreBox (free and paid), an app by Eric Hoffmann that, like Calibre Cloud, accesses Calibre libraries from cloud storage. Unlike Calibre Cloud, it is limited to Dropbox, but CalibreBox supports more than one library at a time, and flexible sorting and filtering. Custom column support for the book detail view, sorting, and filtering by custom columns, and adding more than two libraries are restricted to paid users. The app is built on the design principles of Google's Material Design and is under active development.
Associated apps:
Calibre-go (free), app by Litlcode Studios lets you access your Calibre e-book library from cloud storage and access the library through Calibre-go to browse, sort, search and read books on your mobile. Calibre-go supports multiple libraries across multiple accounts simultaneously.
Calibre Sync (paid), an Android app by BIL Studio that lets you access Calibre libraries from cloud storage (Dropbox, OneDrive, Box, and pCloud), or from SD card. Calibre Sync supports multiple libraries across multiple accounts simultaneously, also allows users to browse, sort, search, filter and download books to read on devices. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Homodes**
Homodes:
Homodes is a genus of moths of the family Erebidae first described by Achille Guenée in 1852.
Taxonomy:
The genus has previously been classified in the subfamily Calpinae of the family Noctuidae.
Description:
Palpi upturned and reaching vertex of head, where the third joint very minute. Antennae ciliated. Thorax smoothly scaled. Abdomen with dorsal tufts on proximal segments. Tibia nearly smooth. Forelegs of male with a tuft of long hair from base of coxa. Forewing with round apex. Hindwings with vein 5 from near center of discocellulars.
Species:
Homodes bracteigutta (Walker, 1862) India, Thailand, Peninsular Malaysia, Borneo, Saleyer, New Guinea, N.Queensland Homodes crocea Guenée, 1852 India, Thailand, Andamans, Sundaland, Sulawesi, Seram, Kei, New Guinea, Bismarcks Homodes fulva Hampson, 1896 Sri Lanka, Borneo Homodes lassula Prout, 1928 Sumatra, Borneo Homodes iomolybda Meyrick, 1889 India (Meghalaya), New Guinea Homodes lithographa Hampson, 1926 Solomon Islands Homodes magnifica Viette, 1958 Madagascar Homodes muluensis Holloway, 2005 Borneo Homodes ornata Roepke, 1938 northern Sulawesi Homodes perilitha Hampson, 1926 southern India, Borneo, Philippines Homodes vivida Guenée, 1852 India, Sri Lanka, Myanmar, Singapore, Borneo, Sulawesi | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**DSDP 367**
DSDP 367:
The DSDP 367 was an area that was drilled as part of the Deep Sea Drilling Project that took place below the Cape Verde Basin.
Location:
The area was drilled from February 22 to March 1, 1975 by the ship Glomar Challenger before DSDP 368 was drilled. Its location was at 12°29.2'N and, 20°02.8'W and is located 370 km southwest of Dakar and 460 km southeast of Praia, Cape Verde and south of the Cape Verde Rise. The seabed is 4,768 meters deep. The drilling carried a total of 984.5 meters of sediment.
Stratigraphy:
At the ocean floor and below consists of several layers including nannomarls (1), clays (2), multicolored silty clay (3), that level located 5,000 meters deep, below are black shales (4a and 4b) and nano-limestones (5a and 5b). Below is the oceanic crust composing basalt (7) just below around 5,800 metres deep.
The top part were formed during the Pleistocene and Miocene age, the second unit were formed during the Late Eocene age, the b subunit were also formed during the Late Paleocene age. The lower units were formed during the Valangian, Oxfordian and Kimmeridgian ages.
Fossil content:
Not including benthic and planktonic (or planctonic) materials. There are types of nanoplanktons (or nanoplanktons) as well as sponge needles.
Foraminifers Planktonic foraminifers are found at around 200 meters below the ocean floor, they include: N 22 – Pleistocene: Globorotalia tumida flexuosa.
N 19 and N 18 – Pliocene (Zanclean): Globigerina rubescens, Globigerinoides conglobatus, Globorotalia crassaformis, Globorotalia digita, Globorotalia exilis, Globorotalia margaritae, Globorotalia miocenica, Globorotalia multicamerata, Globorotalia tosaensis, Globorotalia tumida and Sphaeroidinella dehiscens.
N 12 – Mid to Late Miocene (Messinian/Tortonian): Cassigerinella chipolensis, Globigerina angustiumbilicata and Globigerinoides trilobus.
P 21 – Oligocene: Globigerina ciperoensis, Globigerina ouachitaensis and Globigerina praebulloides.
P 14 – Early to Mid Eocene: Acarinina sp. and Globorotalia subbotinae.
Upper Cretaceous: Gyroidina, Hedbergella amabilis, Hedbergella infracretacea, Hedbergella planispira, Heterohelix and Loeblichella.
Cenomanian: Globigerinelloides caseyi, Guembelitria harrisi, Hedbergella amabilis, Hedbergella brittonensis, Hedbergella delrioensis, Hedbergella globigerinelloides, Hedbergella infracretacea, Hedbergella trochoidea, Heterohelix moremani, Praeglobotruncana and Schackoina cenomana.
Albian: Clavihedbergella simplex, Globigerinelloides, Hedbergella amabilis, Hedbergella delrioensis, Hedbergella infracretacea, Hedbergella planispira, Hedbergella simplicissima and Ticinella primula.
Early Aptian to Barremian: Globigerinelloides, Gubkinella, Hedbergella globigerinelloides, Hedbergella graysonensis, Hedbergella infracretacea and Hedbergella kugleri.
Lower Cretaceous: Dorothia praehauteriviana.
Late Jurassic period: Lenticulina, Nodosaria, Rhabdammina, Spirillina and Spirophthalmidium.
Coccoliths Coccoliths are founded up to 250 meters below the ocean floor, the drilling area, they include: NN 21 and NN 20 – Holocene and Pleistocene: Emiliana huxleyi, Gephyrocapsa oceanica.
NN 19 – Pleistocene: Pseudomiliania lacunosa.
NN 18 and NN 13– Pleocene: Discoaster pentaradiatus, Discoaster surculusaund Discoaster tamalis.
NP 12 –Early Eocene: Coccolithus crassus and Discoaster lodoensis.
Late Cretaceous: Tetralithus obscurus CC 9 – Cenomanian and Late Albian: Chiastozygus amphipons, Corollithion signum, Eiffellithus turriseiffelii, Lithastrinus floralis, Tetralithus obscurus and Vagalapilla matalosa.
CC 8 and CC 7 - Early Albian and Late Aptian: Lithastrinus floralis and Parhabdolithus angustus.
CC 6 – Barremian: Nannoconus colomii.
CC 3 and CC 4 - Hauterivian und Valanginian:Cruciellipsis cuvillieri and Diadorhombus rectus.
CC 1 – Berriasian: Cruciellipsis cuvillieri, Lithraphidites carniolensis and Nannoconus colomii.
NJ 17 – Tithonian: Parhabdolithus embergeri.
NJ 15 a and NJ 15 b – Kimmeridgian and Oxfordian: Callolithus martelae and Cyclagelosphaera margareli.
Radiolaria Several radiolaria were made during the Late Pleistocene, Early Miocene and Early Eocene periods: Late Pleistocene: Axoprunum angelinum, Lamprocyclas maritalis, Ommatartus tetrathalamus, Pterocanium trilobum and Siphocampe corbula.
Late Eocene: Thyrsocyrtis bromia area.
Mid to Late Eocene: Theocampe cryptocephala cryptocephala area and Theocampe mongolifieri rea.
Early Eocene:Buryella clinata area und Phormocyrtis striata striata area.
Early Eocene/Late Pleistocene: Bekoma bidartensis area - Buryella tetradica, Lithocyclia archaea, Lophocyrtis biaurita, Phormocyrtis turgida, Stylosphaera coronate sabaca and Thecosphaerella rotunda.
Albian/Aptian: Lithocampe elegantissima.
Early Cretaceous: Sphaerostylus lanceola.
Berriasian: Sethocapsa trachyostraca-Zone: Sethocapsa cetia and Sethocapsa trachyostraca.
Geological development:
Unlike DSDP 368 which is located 550 km north in the Cape Verde Basin, the Upper Jurassic and the Lower Cretaceous sediments below the black shale of the oceanic crust are founded. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Sequence**
Sequence:
In mathematics, a sequence is an enumerated collection of objects in which repetitions are allowed and order matters. Like a set, it contains members (also called elements, or terms). The number of elements (possibly infinite) is called the length of the sequence. Unlike a set, the same elements can appear multiple times at different positions in a sequence, and unlike a set, the order does matter. Formally, a sequence can be defined as a function from natural numbers (the positions of elements in the sequence) to the elements at each position. The notion of a sequence can be generalized to an indexed family, defined as a function from an arbitrary index set.
Sequence:
For example, (M, A, R, Y) is a sequence of letters with the letter 'M' first and 'Y' last. This sequence differs from (A, R, M, Y). Also, the sequence (1, 1, 2, 3, 5, 8), which contains the number 1 at two different positions, is a valid sequence. Sequences can be finite, as in these examples, or infinite, such as the sequence of all even positive integers (2, 4, 6, ...).
Sequence:
The position of an element in a sequence is its rank or index; it is the natural number for which the element is the image. The first element has index 0 or 1, depending on the context or a specific convention. In mathematical analysis, a sequence is often denoted by letters in the form of an , bn and cn , where the subscript n refers to the nth element of the sequence; for example, the nth element of the Fibonacci sequence F is generally denoted as Fn In computing and computer science, finite sequences are sometimes called strings, words or lists, the different names commonly corresponding to different ways to represent them in computer memory; infinite sequences are called streams. The empty sequence ( ) is included in most notions of sequence, but may be excluded depending on the context.
Examples and notation:
A sequence can be thought of as a list of elements with a particular order. Sequences are useful in a number of mathematical disciplines for studying functions, spaces, and other mathematical structures using the convergence properties of sequences. In particular, sequences are the basis for series, which are important in differential equations and analysis. Sequences are also of interest in their own right, and can be studied as patterns or puzzles, such as in the study of prime numbers.
Examples and notation:
There are a number of ways to denote a sequence, some of which are more useful for specific types of sequences. One way to specify a sequence is to list all its elements. For example, the first four odd numbers form the sequence (1, 3, 5, 7). This notation is used for infinite sequences as well. For instance, the infinite sequence of positive odd integers is written as (1, 3, 5, 7, ...). Because notating sequences with ellipsis leads to ambiguity, listing is most useful for customary infinite sequences which can be easily recognized from their first few elements. Other ways of denoting a sequence are discussed after the examples.
Examples and notation:
Examples The prime numbers are the natural numbers greater than 1 that have no divisors but 1 and themselves. Taking these in their natural order gives the sequence (2, 3, 5, 7, 11, 13, 17, ...). The prime numbers are widely used in mathematics, particularly in number theory where many results related to them exist.
Examples and notation:
The Fibonacci numbers comprise the integer sequence whose elements are the sum of the previous two elements. The first two elements are either 0 and 1 or 1 and 1 so that the sequence is (0, 1, 1, 2, 3, 5, 8, 13, 21, 34, ...).Other examples of sequences include those made up of rational numbers, real numbers and complex numbers. The sequence (.9, .99, .999, .9999, ...), for instance, approaches the number 1. In fact, every real number can be written as the limit of a sequence of rational numbers (e.g. via its decimal expansion). As another example, π is the limit of the sequence (3, 3.1, 3.14, 3.141, 3.1415, ...), which is increasing. A related sequence is the sequence of decimal digits of π, that is, (3, 1, 4, 1, 5, 9, ...). Unlike the preceding sequence, this sequence does not have any pattern that is easily discernible by inspection.
Examples and notation:
Another example of sequences is a sequence of functions, where each member of the sequence is a function whose shape is determined by a natural number indexing that function.
The On-Line Encyclopedia of Integer Sequences comprises a large list of examples of integer sequences.
Examples and notation:
Indexing Other notations can be useful for sequences whose pattern cannot be easily guessed or for sequences that do not have a pattern such as the digits of π. One such notation is to write down a general formula for computing the nth term as a function of n, enclose it in parentheses, and include a subscript indicating the set of values that n can take. For example, in this notation the sequence of even numbers could be written as (2n)n∈N . The sequence of squares could be written as (n2)n∈N . The variable n is called an index, and the set of values that it can take is called the index set.
Examples and notation:
It is often useful to combine this notation with the technique of treating the elements of a sequence as individual variables. This yields expressions like (an)n∈N , which denotes a sequence whose nth element is given by the variable an . For example: st element of nd element rd element th element th element th element ⋮ One can consider multiple sequences at the same time by using different variables; e.g. (bn)n∈N could be a different sequence than (an)n∈N . One can even consider a sequence of sequences: ((am,n)n∈N)m∈N denotes a sequence whose mth term is the sequence (am,n)n∈N An alternative to writing the domain of a sequence in the subscript is to indicate the range of values that the index can take by listing its highest and lowest legal values. For example, the notation 10 denotes the ten-term sequence of squares 100 ) . The limits ∞ and −∞ are allowed, but they do not represent valid values for the index, only the supremum or infimum of such values, respectively. For example, the sequence (an)n=1∞ is the same as the sequence (an)n∈N , and does not contain an additional term "at infinity". The sequence (an)n=−∞∞ is a bi-infinite sequence, and can also be written as (…,a−1,a0,a1,a2,…) In cases where the set of indexing numbers is understood, the subscripts and superscripts are often left off. That is, one simply writes (ak) for an arbitrary sequence. Often, the index k is understood to run from 1 to ∞. However, sequences are frequently indexed starting from zero, as in (ak)k=0∞=(a0,a1,a2,…).
Examples and notation:
In some cases, the elements of the sequence are related naturally to a sequence of integers whose pattern can be easily inferred. In these cases, the index set may be implied by a listing of the first few abstract elements. For instance, the sequence of squares of odd numbers could be denoted in any of the following ways.
Examples and notation:
25 ,…) (a1,a3,a5,…),ak=k2 (a2k−1)k=1∞,ak=k2 (ak)k=1∞,ak=(2k−1)2 ((2k−1)2)k=1∞ Moreover, the subscripts and superscripts could have been left off in the third, fourth, and fifth notations, if the indexing set was understood to be the natural numbers. In the second and third bullets, there is a well-defined sequence (ak)k=1∞ , but it is not the same as the sequence denoted by the expression.
Examples and notation:
Defining a sequence by recursion Sequences whose elements are related to the previous elements in a straightforward way are often defined using recursion. This is in contrast to the definition of sequences of elements as functions of their positions.
To define a sequence by recursion, one needs a rule, called recurrence relation to construct each element in terms of the ones before it. In addition, enough initial elements must be provided so that all subsequent elements of the sequence can be computed by successive applications of the recurrence relation.
The Fibonacci sequence is a simple classical example, defined by the recurrence relation an=an−1+an−2, with initial terms a0=0 and a1=1 . From this, a simple computation shows that the first ten terms of this sequence are 0, 1, 1, 2, 3, 5, 8, 13, 21, and 34.
A complicated example of a sequence defined by a recurrence relation is Recamán's sequence, defined by the recurrence relation if the result is positive and not already in the previous terms, otherwise , with initial term 0.
Examples and notation:
A linear recurrence with constant coefficients is a recurrence relation of the form an=c0+c1an−1+⋯+ckan−k, where c0,…,ck are constants. There is a general method for expressing the general term an of such a sequence as a function of n; see Linear recurrence. In the case of the Fibonacci sequence, one has c0=0,c1=c2=1, and the resulting function of n is given by Binet's formula.
Examples and notation:
A holonomic sequence is a sequence defined by a recurrence relation of the form an=c1an−1+⋯+ckan−k, where c1,…,ck are polynomials in n. For most holonomic sequences, there is no explicit formula for expressing an as a function of n. Nevertheless, holonomic sequences play an important role in various areas of mathematics. For example, many special functions have a Taylor series whose sequence of coefficients is holonomic. The use of the recurrence relation allows a fast computation of values of such special functions.
Examples and notation:
Not all sequences can be specified by a recurrence relation. An example is the sequence of prime numbers in their natural order (2, 3, 5, 7, 11, 13, 17, ...).
Formal definition and basic properties:
There are many different notions of sequences in mathematics, some of which (e.g., exact sequence) are not covered by the definitions and notations introduced below.
Formal definition and basic properties:
Definition In this article, a sequence is formally defined as a function whose domain is an interval of integers. This definition covers several different uses of the word "sequence", including one-sided infinite sequences, bi-infinite sequences, and finite sequences (see below for definitions of these kinds of sequences). However, many authors use a narrower definition by requiring the domain of a sequence to be the set of natural numbers. This narrower definition has the disadvantage that it rules out finite sequences and bi-infinite sequences, both of which are usually called sequences in standard mathematical practice. Another disadvantage is that, if one removes the first terms of a sequence, one needs reindexing the remainder terms for fitting this definition. In some contexts, to shorten exposition, the codomain of the sequence is fixed by context, for example by requiring it to be the set R of real numbers, the set C of complex numbers, or a topological space.Although sequences are a type of function, they are usually distinguished notationally from functions in that the input is written as a subscript rather than in parentheses, that is, an rather than a(n). There are terminological differences as well: the value of a sequence at the lowest input (often 1) is called the "first element" of the sequence, the value at the second smallest input (often 2) is called the "second element", etc. Also, while a function abstracted from its input is usually denoted by a single letter, e.g. f, a sequence abstracted from its input is usually written by a notation such as (an)n∈A , or just as (an).
Formal definition and basic properties:
Here A is the domain, or index set, of the sequence.
Sequences and their limits (see below) are important concepts for studying topological spaces. An important generalization of sequences is the concept of nets. A net is a function from a (possibly uncountable) directed set to a topological space. The notational conventions for sequences normally apply to nets as well.
Finite and infinite The length of a sequence is defined as the number of terms in the sequence.
A sequence of a finite length n is also called an n-tuple. Finite sequences include the empty sequence ( ) that has no elements.
Formal definition and basic properties:
Normally, the term infinite sequence refers to a sequence that is infinite in one direction, and finite in the other—the sequence has a first element, but no final element. Such a sequence is called a singly infinite sequence or a one-sided infinite sequence when disambiguation is necessary. In contrast, a sequence that is infinite in both directions—i.e. that has neither a first nor a final element—is called a bi-infinite sequence, two-way infinite sequence, or doubly infinite sequence. A function from the set Z of all integers into a set, such as for instance the sequence of all even integers ( ..., −4, −2, 0, 2, 4, 6, 8, ... ), is bi-infinite. This sequence could be denoted (2n)n=−∞∞ Increasing and decreasing A sequence is said to be monotonically increasing if each term is greater than or equal to the one before it. For example, the sequence (an)n=1∞ is monotonically increasing if and only if an+1 ≥ an for all n ∈ N. If each consecutive term is strictly greater than (>) the previous term then the sequence is called strictly monotonically increasing. A sequence is monotonically decreasing if each consecutive term is less than or equal to the previous one, and is strictly monotonically decreasing if each is strictly less than the previous. If a sequence is either increasing or decreasing it is called a monotone sequence. This is a special case of the more general notion of a monotonic function.
Formal definition and basic properties:
The terms nondecreasing and nonincreasing are often used in place of increasing and decreasing in order to avoid any possible confusion with strictly increasing and strictly decreasing, respectively.
Formal definition and basic properties:
Bounded If the sequence of real numbers (an) is such that all the terms are less than some real number M, then the sequence is said to be bounded from above. In other words, this means that there exists M such that for all n, an ≤ M. Any such M is called an upper bound. Likewise, if, for some real m, an ≥ m for all n greater than some N, then the sequence is bounded from below and any such m is called a lower bound. If a sequence is both bounded from above and bounded from below, then the sequence is said to be bounded.
Formal definition and basic properties:
Subsequences A subsequence of a given sequence is a sequence formed from the given sequence by deleting some of the elements without disturbing the relative positions of the remaining elements. For instance, the sequence of positive even integers (2, 4, 6, ...) is a subsequence of the positive integers (1, 2, 3, ...). The positions of some elements change when other elements are deleted. However, the relative positions are preserved.
Formal definition and basic properties:
Formally, a subsequence of the sequence (an)n∈N is any sequence of the form (ank)k∈N , where (nk)k∈N is a strictly increasing sequence of positive integers.
Other types of sequences Some other types of sequences that are easy to define include: An integer sequence is a sequence whose terms are integers.
A polynomial sequence is a sequence whose terms are polynomials.
A positive integer sequence is sometimes called multiplicative, if anm = an am for all pairs n, m such that n and m are coprime. In other instances, sequences are often called multiplicative, if an = na1 for all n. Moreover, a multiplicative Fibonacci sequence satisfies the recursion relation an = an−1 an−2.
A binary sequence is a sequence whose terms have one of two discrete values, e.g. base 2 values (0,1,1,0, ...), a series of coin tosses (Heads/Tails) H,T,H,H,T, ..., the answers to a set of True or False questions (T, F, T, T, ...), and so on.
Limits and convergence:
An important property of a sequence is convergence. If a sequence converges, it converges to a particular value known as the limit. If a sequence converges to some limit, then it is convergent. A sequence that does not converge is divergent.
Limits and convergence:
Informally, a sequence has a limit if the elements of the sequence become closer and closer to some value L (called the limit of the sequence), and they become and remain arbitrarily close to L , meaning that given a real number d greater than zero, all but a finite number of the elements of the sequence have a distance from L less than d For example, the sequence {\textstyle a_{n}={\frac {n+1}{2n^{2}}}} shown to the right converges to the value 0. On the other hand, the sequences {\textstyle b_{n}=n^{3}} (which begins 1, 8, 27, …) and cn=(−1)n (which begins −1, 1, −1, 1, …) are both divergent.
Limits and convergence:
If a sequence converges, then the value it converges to is unique. This value is called the limit of the sequence. The limit of a convergent sequence (an) is normally denoted lim {\textstyle \lim _{n\to \infty }a_{n}} . If (an) is a divergent sequence, then the expression lim {\textstyle \lim _{n\to \infty }a_{n}} is meaningless.
Formal definition of convergence A sequence of real numbers (an) converges to a real number L if, for all ε>0 , there exists a natural number N such that for all n≥N we have |an−L|<ε.
Limits and convergence:
If (an) is a sequence of complex numbers rather than a sequence of real numbers, this last formula can still be used to define convergence, with the provision that |⋅| denotes the complex modulus, i.e. |z|=z∗z . If (an) is a sequence of points in a metric space, then the formula can be used to define convergence, if the expression |an−L| is replaced by the expression dist (an,L) , which denotes the distance between an and L Applications and important results If (an) and (bn) are convergent sequences, then the following limits exist, and can be computed as follows: lim lim lim n→∞bn lim lim n→∞an for all real numbers c lim lim lim n→∞bn) lim lim lim n→∞bn , provided that lim n→∞bn≠0 lim lim n→∞an)p for all p>0 and an>0 Moreover: If an≤bn for all n greater than some N , then lim lim n→∞bn (Squeeze Theorem)If (cn) is a sequence such that an≤cn≤bn for all n>N and lim lim n→∞bn=L ,then (cn) is convergent, and lim n→∞cn=L If a sequence is bounded and monotonic then it is convergent.
Limits and convergence:
A sequence is convergent if and only if all of its subsequences are convergent.
Limits and convergence:
Cauchy sequences A Cauchy sequence is a sequence whose terms become arbitrarily close together as n gets very large. The notion of a Cauchy sequence is important in the study of sequences in metric spaces, and, in particular, in real analysis. One particularly important result in real analysis is Cauchy characterization of convergence for sequences: A sequence of real numbers is convergent (in the reals) if and only if it is Cauchy.In contrast, there are Cauchy sequences of rational numbers that are not convergent in the rationals, e.g. the sequence defined by x1 = 1 and xn+1 = xn + 2/xn/2 is Cauchy, but has no rational limit, cf. here. More generally, any sequence of rational numbers that converges to an irrational number is Cauchy, but not convergent when interpreted as a sequence in the set of rational numbers.
Limits and convergence:
Metric spaces that satisfy the Cauchy characterization of convergence for sequences are called complete metric spaces and are particularly nice for analysis.
Infinite limits In calculus, it is common to define notation for sequences which do not converge in the sense discussed above, but which instead become and remain arbitrarily large, or become and remain arbitrarily negative. If an becomes arbitrarily large as n→∞ , we write lim n→∞an=∞.
In this case we say that the sequence diverges, or that it converges to infinity. An example of such a sequence is an = n.
If an becomes arbitrarily negative (i.e. negative and large in magnitude) as n→∞ , we write lim n→∞an=−∞ and say that the sequence diverges or converges to negative infinity.
Series:
A series is, informally speaking, the sum of the terms of a sequence. That is, it is an expression of the form {\textstyle \sum _{n=1}^{\infty }a_{n}} or a1+a2+⋯ , where (an) is a sequence of real or complex numbers. The partial sums of a series are the expressions resulting from replacing the infinity symbol with a finite number, i.e. the Nth partial sum of the series {\textstyle \sum _{n=1}^{\infty }a_{n}} is the number SN=∑n=1Nan=a1+a2+⋯+aN.
Series:
The partial sums themselves form a sequence (SN)N∈N , which is called the sequence of partial sums of the series {\textstyle \sum _{n=1}^{\infty }a_{n}} . If the sequence of partial sums converges, then we say that the series {\textstyle \sum _{n=1}^{\infty }a_{n}} is convergent, and the limit lim {\textstyle \lim _{N\to \infty }S_{N}} is called the value of the series. The same notation is used to denote a series and its value, i.e. we write lim {\textstyle \sum _{n=1}^{\infty }a_{n}=\lim _{N\to \infty }S_{N}}
Use in other fields of mathematics:
Topology Sequences play an important role in topology, especially in the study of metric spaces. For instance: A metric space is compact exactly when it is sequentially compact.
A function from a metric space to another metric space is continuous exactly when it takes convergent sequences to convergent sequences.
A metric space is a connected space if and only if, whenever the space is partitioned into two sets, one of the two sets contains a sequence converging to a point in the other set.
A topological space is separable exactly when there is a dense sequence of points.Sequences can be generalized to nets or filters. These generalizations allow one to extend some of the above theorems to spaces without metrics.
Product topology The topological product of a sequence of topological spaces is the cartesian product of those spaces, equipped with a natural topology called the product topology.
Use in other fields of mathematics:
More formally, given a sequence of spaces (Xi)i∈N , the product space := ∏i∈NXi, is defined as the set of all sequences (xi)i∈N such that for each i, xi is an element of Xi . The canonical projections are the maps pi : X → Xi defined by the equation pi((xj)j∈N)=xi . Then the product topology on X is defined to be the coarsest topology (i.e. the topology with the fewest open sets) for which all the projections pi are continuous. The product topology is sometimes called the Tychonoff topology.
Use in other fields of mathematics:
Analysis In analysis, when talking about sequences, one will generally consider sequences of the form or (x0,x1,x2,…) which is to say, infinite sequences of elements indexed by natural numbers.
Use in other fields of mathematics:
A sequence may start with an index different from 1 or 0. For example, the sequence defined by xn = 1/log(n) would be defined only for n ≥ 2. When talking about such infinite sequences, it is usually sufficient (and does not change much for most considerations) to assume that the members of the sequence are defined at least for all indices large enough, that is, greater than some given N.
Use in other fields of mathematics:
The most elementary type of sequences are numerical ones, that is, sequences of real or complex numbers. This type can be generalized to sequences of elements of some vector space. In analysis, the vector spaces considered are often function spaces. Even more generally, one can study sequences with elements in some topological space.
Use in other fields of mathematics:
Sequence spaces A sequence space is a vector space whose elements are infinite sequences of real or complex numbers. Equivalently, it is a function space whose elements are functions from the natural numbers to the field K, where K is either the field of real numbers or the field of complex numbers. The set of all such functions is naturally identified with the set of all possible infinite sequences with elements in K, and can be turned into a vector space under the operations of pointwise addition of functions and pointwise scalar multiplication. All sequence spaces are linear subspaces of this space. Sequence spaces are typically equipped with a norm, or at least the structure of a topological vector space.
Use in other fields of mathematics:
The most important sequences spaces in analysis are the ℓp spaces, consisting of the p-power summable sequences, with the p-norm. These are special cases of Lp spaces for the counting measure on the set of natural numbers. Other important classes of sequences like convergent sequences or null sequences form sequence spaces, respectively denoted c and c0, with the sup norm. Any sequence space can also be equipped with the topology of pointwise convergence, under which it becomes a special kind of Fréchet space called an FK-space.
Use in other fields of mathematics:
Linear algebra Sequences over a field may also be viewed as vectors in a vector space. Specifically, the set of F-valued sequences (where F is a field) is a function space (in fact, a product space) of F-valued functions over the set of natural numbers.
Abstract algebra Abstract algebra employs several types of sequences, including sequences of mathematical objects such as groups or rings.
Free monoid If A is a set, the free monoid over A (denoted A*, also called Kleene star of A) is a monoid containing all the finite sequences (or strings) of zero or more elements of A, with the binary operation of concatenation. The free semigroup A+ is the subsemigroup of A* containing all elements except the empty sequence.
Exact sequences In the context of group theory, a sequence G0→f1G1→f2G2→f3⋯→fnGn of groups and group homomorphisms is called exact, if the image (or range) of each homomorphism is equal to the kernel of the next: im(fk)=ker(fk+1) The sequence of groups and homomorphisms may be either finite or infinite.
A similar definition can be made for certain other algebraic structures. For example, one could have an exact sequence of vector spaces and linear maps, or of modules and module homomorphisms.
Spectral sequences In homological algebra and algebraic topology, a spectral sequence is a means of computing homology groups by taking successive approximations. Spectral sequences are a generalization of exact sequences, and since their introduction by Jean Leray (1946), they have become an important research tool, particularly in homotopy theory.
Set theory An ordinal-indexed sequence is a generalization of a sequence. If α is a limit ordinal and X is a set, an α-indexed sequence of elements of X is a function from α to X. In this terminology an ω-indexed sequence is an ordinary sequence.
Computing In computer science, finite sequences are called lists. Potentially infinite sequences are called streams. Finite sequences of characters or digits are called strings.
Use in other fields of mathematics:
Streams Infinite sequences of digits (or characters) drawn from a finite alphabet are of particular interest in theoretical computer science. They are often referred to simply as sequences or streams, as opposed to finite strings. Infinite binary sequences, for instance, are infinite sequences of bits (characters drawn from the alphabet {0, 1}). The set C = {0, 1}∞ of all infinite binary sequences is sometimes called the Cantor space.
Use in other fields of mathematics:
An infinite binary sequence can represent a formal language (a set of strings) by setting the n th bit of the sequence to 1 if and only if the n th string (in shortlex order) is in the language. This representation is useful in the diagonalization method for proofs. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Quixotism**
Quixotism:
Quixotism ( or ; adj. quixotic) is impracticality in pursuit of ideals, especially those ideals manifested by rash, lofty and romantic ideas or extravagantly chivalrous action. It also serves to describe an idealism without regard to practicality. An impulsive person or act might be regarded as quixotic.
Quixotism is usually related to "over-idealism", meaning an idealism that doesn't take consequence or absurdity into account. It is also related to naïve romanticism and to utopianism.
Origin:
Quixotism as a term or a quality appeared after the publication of Don Quixote in 1605. Don Quixote, the hero of this novel, written by Spanish author Miguel de Cervantes Saavedra, dreams up a romantic ideal world which he believes to be real, and acts on this idealism, which most famously leads him into imaginary fights with windmills that he regards as giants, leading to the related metaphor of "tilting at windmills".
Origin:
Already in the 17th century the term quixote was used to describe a person who does not distinguish between reality and imagination. The poet John Cleveland wrote in 1644, in his book The character of a London diurnall: The Quixotes of this Age fight with the Wind-mills of their own Heads.
The word quixotism is mentioned, for the first time, in Pulpit Popery, True Popery (1688): ...all the Heroical Fictions of Ecclesiastical Quixotism...
Spanish language opposes quijotesco ("Quixotic") with sanchopancesco ("lacking idealism, accommodating and chuckling" after Sancho Panza). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Zinc pest**
Zinc pest:
Zinc pest (from German Zinkpest), also known as zinc rot and zamak rot, is a destructive, intercrystalline corrosion process of zinc alloys containing lead impurities. While impurities of the alloy are the primary cause of the problem, environmental conditions such as high humidity (greater than 65%) may accelerate the process.It was first discovered to be a problem in 1923, and primarily affects die-cast zinc articles that were manufactured during the 1920s through 1950s. The New Jersey Zinc Company developed zamak alloys in 1929 using 99.99% pure zinc metal to avoid the problem, and articles made after 1960 are usually considered free of the risk of zinc pest since the use of purer materials and more controlled manufacturing conditions make zinc pest degradation unlikely.Affected objects may show surface irregularities such as small cracks and fractures, blisters or pitting. Over time, the material slowly expands, cracking, buckling and warping in an irreversible process that makes the object exceedingly brittle and prone to fracture, and can eventually shatter the object, destroying it altogether. Due to the expansion process, attached normal material may also be damaged. The occurrence and severity of zinc pest in articles made of susceptible zinc alloys depends both on the concentration of lead impurities in the metal and on the storage conditions of the article in the ensuing decades. Zinc pest is dreaded by collectors of vintage die-cast model trains, toys, or radios, because rare or otherwise valuable items can inescapably be rendered worthless as the process of zinc pest destroys them. Because castings of the same object were usually made from various batches of metal over the production process, some examples of a given toy or model may survive today completely unaffected, while other identical examples may have completely disintegrated. It has also affected carburetors, hubcaps, door handles and automobile trim on cars of the 1920s and 1930s.
Zinc pest:
Since the 1940s, some model railroad hobbyists have claimed, with varying degrees of success, that a method of "pickling" zinc alloy parts by soaking them in vinegar or oxalic acid solution for several minutes before painting and assembling them could prevent or delay the effects of zinc pest. Engine parts of older vehicles or airplanes, and military medals made of zinc alloys, may also be affected. In addition, the post-1982 copper-plated zinc Lincoln cents have been known to be affected.
Zinc pest:
Zinc pest is not related to tin pest, and is also different from a superficial white corrosion oxidation process ("Weissrost") that affects some zinc articles. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Vaneless ion wind generator**
Vaneless ion wind generator:
A vaneless ion wind generator or power fence is a device that generates electrical energy by using the wind to move charged particles across an electric field.
Vaneless ion wind generator:
Ion wind generators are not commercially available, though working prototypes and proofs of concept have been created. Several prototypes exist in the Netherlands, one of which resides in Delft University of Technology, whose researchers developed some of the underlying technology. Ion wind generators are currently experimental, while conventional wind turbines are the most common form of wind energy generation. But ion wind generators, which have no moving parts, could be used in urban settings where wind turbines are impractical due to vibrational noise, moving shadows, and danger posed to birds.
History:
Lord Kelvin's Thunderstorm One of the earliest examples of electrostatic energy generation is found in Lord Kelvin's Thunderstorm, a device invented in 1867. Similar to ion wind generators, the Thunderstorm used water to carry charges and generate energy through related principles. However, the Thunderstorm relied on the force of gravity and two oppositely charged reservoirs to generate a voltage difference. Though they are not identical in operation, Lord Kelvin's Thunderstorm demonstrates the behavior of water and concepts of electrostatics that underpin modern ion wind generators.
Design and construction:
Theoretical operation Ion wind generators use the force of the wind to move charged particles, typically water, against the force of an electric field. This increases the potential energy of the particles, which can be likened to moving a mass upwards against the force of gravity. The method of collecting the energy varies by implementation. The design of ion wind generators eliminates the intermediate conversion of mechanical energy undergone in wind turbines. Wind turbines use the kinetic energy of the wind to rotate several blades about a rotor. The rotor's mechanical energy is converted into electrical energy by an electric generator.
Design and construction:
Conversion between different forms of energy necessitates some energy loss, either to the environment or in a useless form, and fewer conversions improve theoretical output.
Design and construction:
Simplified analytical model Researchers from Delft University of Technology devised an equation to model the behavior of the water droplets as they move through the air in order to optimize the system mathematically and run computer simulations. For the purposes of the model, a simple electrode configuration and uniform electric field is assumed, wherein the electric force exerted on the particles will be directly opposite that of the wind.
Design and construction:
Each particle is acted upon by the force of gravity, where mi is the mass of the ith droplet and g→ is the gravitational acceleration of Earth. The model assumes mi is constant and does not take evaporation into account. The atmosphere also exerts a force in the form of buoyancy as the droplets fall, where Vd is the volume of the droplet and ρa is the air density. The droplets are also acted upon by the wind, where Cd is the drag coefficient, vw is the wind speed, and vd is the droplet speed. The equation may be simplified in cases of laminar flow, which can be expressed using the Reynolds Number (Re), which is used in fluid mechanics to determine flow patterns. The flow is considered laminar when the Reynolds Number is less than 1, where ηa is the viscosity of air. When the flow is indeed laminar, the drag force can be calculated using Stoke's law, where Cc is the Cunningham slip correction factor, which is assumed to be 1 for particles greater than 1μm in diameter.
Design and construction:
The electric force acting on the droplets is affected by both the external electric field ( Eext ) of the device's electrodes, where qi is the charge of the ith droplet, and the electric fields of other charged droplets, where ri,j is the distance between droplet i and droplet j. The sum of these forces represents the researchers' full equation, where Fi is the total force exerted on the ith droplet and ai is the acceleration of the ith droplet. The work done on the ith droplet can be calculated using the previous equation, where dl is the droplet's displacement. The researchers use this to calculate the potential energy difference for the droplet. The sum of work done on each droplet yields the total energy generated from the wind.
Design and construction:
Implementations There are two mainstream implementations of ion wind generators. The first, patented by Alvin Marks in 1977, was a twofold device comprising a charging system and separate collector. The EWICON is a derivative of the design that allows the system to function without the need for a separate collector.
Alvin Marks' patent A grounded charging system produces a cloud of charged particles. The wind carries the particles toward a conducting collector. The collector is insulated by its non-conducting mechanical support. Though the collector is initially neutral, the particles transfer their charge upon contact, increasing the collector's potential energy.
Design and construction:
The charged particles and the collector, now also charged, form an electric field which exerts a force on the particles in the opposite direction of the wind. Though the force of the wind initially exceeds the force of the electric field, the continuous flow of particles increases the force of the electric field. The force may become strong enough to move the particles back towards the charging system, or they may simply pass by the collector. The particles which never reach the collector do not contribute to the net energy generation.
Design and construction:
The system performs at maximum efficiency when all particles reach the collector. Adjusting variables such as wind speed and collector size can improve the performance of the system.
Design and construction:
EWICON (Electrostatic Wind Energy Converter) The EWICON functions using the same principles as the previous implementation, but abandons the collector. Instead, the EWICON is insulated from the Earth, and releases charged particles into the air. The dispersal of negatively charged particles from an initially neutral system increases its potential energy. Once the charging system has a polarity, which is opposite to that of the particles, an attractive force is exerted. If there is little wind, the force may transport the particles back to the charging system, losing the net energy gained from their dispersal.
Design and construction:
The EWICON system performs at maximum efficiency when all particles leave the charging system and reach the Earth, which acts as the collector in lieu of a secondary system.A group of researchers from Delft University of Technology devised the system. One prototype of the device was installed on the university campus, and two more sit atop the Stadstimmerhuis 010 building located in Rotterdam. The prototypes were designed by Mecanoo, a local architecture firm in Delft.
Design and construction:
The Dutch Windwheel The Dutch Windwheel is a building design that is expected to incorporate EWICON technology. The plans were proposed by a partnership of three Rotterdam companies through the Dutch Windwheel Corp., who expected the building to be completed by 2022, but has not begun construction. The structure is intended to display multiple environmentally-friendly technologies, including rainwater capture, wetland water filtration, and solar energy. The center of the circular building is reserved for wind power generation through the use of a large-scale ion wind generator based on the EWICON implementation. The efficiency and power generation of the system at such a scale is not known, but the Dutch Windwheel Corp. expects the building to generate more energy than it consumes.
Comparison with wind turbines:
Ion wind generators and wind turbines share some of the same advantages and disadvantages. Both are subject to the conditions of the wind, and are unable to generate electricity if the weather conditions are not favorable. This can be mitigated to some degree with strategic placement of the devices in areas with more consistent wind speed.
Comparison with wind turbines:
Advantages Ion wind generators are typically much smaller than wind turbines. Many wind turbine models exceed 400 feet (122 m) in height. Their size and complexity lead to high maintenance costs, which, when combined with the cost of operation, may account for a quarter of the total cost per kilowatt-hour. Wind turbines also produce noise which may disturb residents in the vicinity. The aerodynamic properties of wind turbine blades and inner mechanical workings produce the noise, yet both features are not present in ion wind generators. Quieter operation has led researchers to consider using the technology in urban environments. The bladeless design of ion wind generators could make wind power more environmentally friendly, as current "wind power plants represent a risk of bird mortality." Wind turbines have maximum speeds of operation which vary by design. Wind turbines shut down when "cut-out" speeds are exceeded to prevent damage. Therefore turbines are unable to generate energy in high speed winds which fall beyond the window of performance, while ion wind generators can theoretically continue to operate.
Comparison with wind turbines:
Disadvantages The technology is still in its nascence, and ion wind generators are not as efficient as conventional wind turbines. During tests conducted in 2005, the EWICON was unable to match wind turbine output. Researchers were able to demonstrate "a conversion of 7% of the wind energy into electrical energy, whereas conventional wind turbine systems have an efficiency of 45% at their rated speeds. Improvements are suggested that could lead to an efficiency of the EWICON in the range of 25–30%." At the 2005 International Conference on Future Power Systems, suggestions for future advancements included changes to the method of electrohydrodynamic atomization, or electrospray, and designing a more dense array of nozzles. Tests have yet to indicate that the technology has developed enough to rival wind turbines in efficiency. Several prototypes have been built for testing and experimentation, but researchers hope to build a larger device with greater power output. While the current level of development does not surpass wind turbines in efficiency, the technology could contribute to the energy mix in urban environments where a wind turbine may be impractical. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Dibunate**
Dibunate:
Dibunate is a cough suppressant. As the sodium salt, it has been marketed under the name Becantyl (in the United Kingdom), Becantex (in continental Europe), or Linctussal with a dosage of 20 to 30 mg, as either syrup or tablets.Similar to benzonatate, it is a peripherally acting drug. It has not been reported to cause sedation, euphoria, habituation, or respiratory depression, unlike narcotic antitussives such as codeine. It may work by blocking afferent signals in the reflex arc which controls cough. Nausea is rarely seen as an adverse effect. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Autocorrelation (words)**
Autocorrelation (words):
In combinatorics, a branch of mathematics, the autocorrelation of a word is the set of periods of this word. More precisely, it is a sequence of values which indicate how much the end of a word looks likes the beginning of a word. This value can be used to compute, for example, the average value of the first occurrence of this word in a random string.
Definition:
In this article, A is an alphabet, and w=w1…wn a word on A of length n. The autocorrelation of w can be defined as the correlation of w with itself. However, we redefine this notion below.
Definition:
Autocorrelation vector The autocorrelation vector of w is c=(c0,…,cn−1) , with ci being 1 if the prefix of length n−i equals the suffix of length n−i , and with ci being 0 otherwise. That is ci indicates whether wi+1…wn=w1…wn−i For example, the autocorrelation vector of aaa is (1,1,1) since, clearly, for i being 0, 1 or 2, the prefix of length n−i is equal to the suffix of length n−i . The autocorrelation vector of abb is (1,0,0) since no strict prefix is equal to a strict suffix. Finally, the autocorrelation vector of aabbaa is 100011, as shown in the following table: Note that c0 is always equal to 1, since the prefix and the suffix of length n are both equal to the word w . Similarly, cn−1 is 1 if and only if the first and the last letters are the same.
Definition:
Autocorrelation polynomial The autocorrelation polynomial of w is defined as c(z)=c0z0+⋯+cn−1zn−1 . It is a polynomial of degree at most n−1 . For example, the autocorrelation polynomial of aaa is 1+z+z2 and the autocorrelation polynomial of abb is 1 . Finally, the autocorrelation polynomial of aabbaa is 1+z4+z5
Property:
We now indicate some properties which can be computed using the autocorrelation polynomial.
Property:
First occurrence of a word in a random string Suppose that you choose an infinite sequence s of letters of A , randomly, each letter with probability 1|A| , where |A| is the number of letters of A . Let us call E the expectation of the first occurrence of ? m in s ? . Then E equals |A|nc(1|A|) . That is, each subword v of w which is both a prefix and a suffix causes the average value of the first occurrence of w to occur |A||v| letters later. Here |v| is the length of v.
Property:
For example, over the binary alphabet A={a,b} , the first occurrence of aa is at position 22(1+12)=6 while the average first occurrence of ab is at position 22(1)=4 . Intuitively, the fact that the first occurrence of aa is later than the first occurrence of ab can be explained in two ways: We can consider, for each position p , what are the requirement for w 's first occurrence to be at p The first occurrence of ab can be at position 1 in only one way in both case. If s starts with w . This has probability 14 for both considered values of w The first occurrence of ab is at position 2 if the prefix of s of length 3 is aab or is bab . However, the first occurrence of aa is at position 2 if and only if the prefix of s of length 3 is baa . (Note that the first occurrence of aa in aaa is at position 1.).
Property:
In general, the number of prefixes of length n+1 such that the first occurrence of aa is at position n is smaller for aa than for ba . This explain why, on average, the first aa arrive later than the first ab We can also consider the fact that the average number of occurrences of w in a random string of length l is |A|l−n . This number is independent of the autocorrelation polynomial. An occurrence of w may overlap another occurrence in different ways. More precisely, each 1 in its autocorrelation vector correspond to a way for occurrence to overlap. Since many occurrences of w can be packed together, using overlapping, but the average number of occurrences does not change, it follows that the distance between two non-overlapping occurrences is greater when the autocorrelation vector contains many 1's.
Property:
Ordinary generating functions Autocorrelation polynomials allows to give simple equations for the ordinary generating functions (OGF) of many natural questions.
The OGF of the languages of words not containing w is c(z)zn+(1−|A|z)c(z) The OGF of the languages of words containing w is zn(1−|A|z)(zn+(1−|A|z)c(z)) The OGF of the languages of words containing a single occurrence of w , at the end of the word is znzn+(1−|A|z)c(z) | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Ameloblastic fibro-odontoma**
Ameloblastic fibro-odontoma:
The ameloblastic fibro-odontoma (AFO) is essentially a benign tumor with the features characteristic of ameloblastic fibroma along with enamel and dentin (hard tissues). Though it is generally regarded as benign, there have been cases of its malignant transformation into ameloblastic fibrosarcoma and odontogenic sarcoma. Cahn LR and Blum T, believed in "maturation theory", which suggested that AFO was an intermediate stage and eventually developed during the period of tooth formation to a complex odontoma thus, being a hamartoma.World Health Organization (WHO) defines AFO as a neoplasm consisting of odontogenic ectomesenchyme resembling the dental papilla, epithelial strands and nest resembling dental lamina and enamel organ conjunction with the presence of dentine and enamel. There is a consensus that AFO should be grouped under Odontomas. This is because once the hard tissues start forming it will eventually lead to formation of Odontomas. The Recent WHO classification published in 2017 has grouped AFDs into odontomes. According to Tekkesin S et al, combination of age and lesion size should be used to distinguish between lesions of a true neoplastic nature and hamartomatous formation.
History:
Initially, AFO was called as ameloblastic odontoma. Hooker in 1967 first used the term ameloblastic fibro-odontoma. WHO classified odontogenic tumors for the first time in 1971 (1stedition). In 2005, mixed tumors were included in the classification of odontogenic tumors by WHO as tumors having odontogenic epithelium along with odontogenic ectomesenchyme. This includes ameloblastic fibroma (AF), ameloblastic fibro-dentinoma (AFD), ameloblastic fibro-odontoma (AFO) and odontoma amongst others. According to the continuum concept, AF, AFD, AFO and odontoma lie in the same spectrum but at different sides depending upon the maturation with AF and odontoma lying completely opposite whereas AFD and AFO lying between them. AFO consists of odontogenic epithelium along with enamel and dentin. Such inductive changes along with proliferating odontogenic epithelium warrant AFO to be regarded as a separate entity.
Histogenesis:
AFO is regarded as a benign mixed odontogenic tumor. Cahn and Blum first advocated continuum concept but this was later rejected. AFO is now considered as a distinct entity. AFO exhibits the same benign biologic behavior as that of AF, along with inductive changes that lead to the formation of both dentin and enamel.
Prevalence:
AFO is a rare odontogenic tumor and accounts for around 2% of all the jaw tumors. This neoplasm usually occurs in the first and second decade of life and is most common in the posterior region of mandible. There appears no gender predilection. Slootweg PJ established that the data on age, site, and sex were consistent with the concept that the AFO was an immature complex odontoma, thereby indicating that AFO was a hamartoma.
Clinical presentation:
AFO is seen in younger age groups, commonly in the second decade of life. Failure and failure of tooth eruption are the most common presenting complaints. In case of large swellings, it may show deformity and show displacement of erupted teeth. Pain and paresthesia are not features of AFO. Accidental discovery during radiographical investigation may be possible. The lesion occurs most commonly in the posterior region of the jaw (mandible). Massive maxillary AFOs cause destruction of the sinus, facial disfigurement, perforated the cortical plates or extended to the orbital floor-pterygoid region.
Radiographical features:
AFO is essentially a radiolucent lesion that is generally unilocular and rarely multilocular. The radiolucency contains calcified material depicting as radio-opaque areas of variable size. The amount of calcified material present is variable and hence the lesion may be either radiolucent, radio-opaque or mixed. Associated crown of the unerupted tooth may be evident.
Histopathology:
The cell-rich mesenchymal tissue resembles the primitive dental papilla. The odontogenic epithelium may be either arranged as follicles that resemble developing enamel organ or they may be seen as cords or strands. The peripheral cells of the follicles resemble ameloblast like cells The odontogenic epithelium is scattered in a loose connective tissue stroma that closely resembles dental papilla. Calcified material may resemble enamel or dentin matrix. The mature lesions may show enamel or dentin aggregates. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Starter (engine)**
Starter (engine):
A starter (also self-starter, cranking motor, or starter motor) is a device used to rotate (crank) an internal-combustion engine so as to initiate the engine's operation under its own power. Starters can be electric, pneumatic, or hydraulic. The starter can also be another internal-combustion engine in the case, for instance, of very large engines, or diesel engines in agricultural or excavation applications.Internal combustion engines are feedback systems, which, once started, rely on the inertia from each cycle to initiate the next cycle. In a four-stroke engine, the third stroke releases energy from the fuel, powering the fourth (exhaust) stroke and also the first two (intake, compression) strokes of the next cycle, as well as powering the engine's external load. To start the first cycle at the beginning of any particular session, the first two strokes must be powered in some other way than from the engine itself. The starter motor is used for this purpose and it is not required once the engine starts running and its feedback loop becomes self-sustaining.
History:
Before the advent of the starter motor, engines were started by various methods including wind-up springs, gunpowder cylinders, and human-powered techniques such as a removable crank handle which engaged the front of the crankshaft, pulling on an airplane propeller, or pulling a cord that was wound around an open-face pulley.
History:
The hand-crank method was commonly used to start engines, but it was inconvenient, difficult, and dangerous. The behavior of an engine during starting is not always predictable. The engine can kick back, causing sudden reverse rotation. Many manual starters included a one-directional slip or release provision so that once engine rotation began, the starter would disengage from the engine. In the event of a kickback, the reverse rotation of the engine could suddenly engage the starter, causing the crank to unexpectedly and violently jerk, possibly injuring the operator. For cord-wound starters, a kickback could pull the operator towards the engine or machine, or swing the starter cord and handle at high speed around the starter pulley. Even though cranks had an overrun mechanism, when the engine started, the crank could begin to spin along with the crankshaft and potentially strike the person cranking the engine. Additionally, care had to be taken to retard the spark in order to prevent backfiring; with an advanced spark setting, the engine could kick back (run in reverse), pulling the crank with it, because the overrun safety mechanism works in one direction only.
History:
Although users were advised to cup their fingers and thumb under the crank and pull up, it felt natural for operators to grasp the handle with the fingers on one side, the thumb on the other. Even a simple backfire could result in a broken thumb; it was possible to end up with a broken wrist, a dislocated shoulder or worse. Moreover, increasingly larger engines with higher compression ratios made hand cranking a more physically demanding endeavour.
History:
The first electric starter was installed on an Arnold, an adaptation of the Benz Velo, built in 1896 in East Peckham, England, by electrical engineer H. J. Dowsing.In 1903, Clyde J. Coleman invented and patented the first electric starter in America U.S. Patent 0,745,157.In 1911, Charles F. Kettering, with Henry M. Leland, of Dayton Engineering Laboratories Company (DELCO), invented and filed U.S. Patent 1,150,523 for an electric starter in America. (Kettering had replaced the hand crank on NCR's cash registers with an electric motor five years earlier.) One aspect of the invention lay in the realization that a relatively small motor, driven with higher voltage and current than would be feasible for continuous operation, could deliver enough power to crank the engine for starting. At the voltage and current levels required, such a motor would burn out in a few minutes of continuous operation, but not during the few seconds needed to start the engine. The starters were first installed on the Cadillac Model Thirty in 1912, with the same system being adopted by Lanchester later that year. These starters also worked as generators once the engine was running, a concept that is now being revived in hybrid vehicles.
History:
Although the electric starter motor was to come to dominate the car market, in 1912, there were several competing types of starter, with the Adams, S.C.A.T. and Wolseley cars having direct air starters, and Sunbeam introducing an air starter motor with similar approach to that used for the Delco and Scott-Crossley electrical starter motors (i.e. engaging with a toothed ring on the flywheel). The Star and Adler cars had spring motors (sometimes referred to as clockwork motors), which used the energy stored in a spring driving through a reduction gear. If the car failed to start, the starter handle could be used to wind up the spring for a further attempt.
History:
One of the innovations on the first Dodge car, the Model 30-35 at its introduction in 1914 was an electric starter and electric lighting with a 12-volt system (against the six volts that was usual at the time) as a standard fitment on what was a relatively low-priced car. The Dodge used a combined starter-generator unit, with a direct current dynamo permanently coupled by gears to the engine's crankshaft. A system of electrical relays allowed this to be driven as a motor to rotate the engine for starting, and once the starter button was released the controlling switchgear returned the unit to operation as a generator. Because the starter-generator was directly coupled to the engine it did not need a method of engaging and disengaging the motor drive. It thus suffered negligible mechanical wear and was virtually silent in operation. The starter-generator remained a feature of Dodge cars until 1929. The disadvantage of the design was that, as a dual-purpose device, the unit was limited in both its power as a motor and its output as a generator, which became a problem as engine size and electrical demands on cars increased. Controlling the switch between motor and generator modes required dedicated and relatively complex switchgear which was more prone to failure than the heavy-duty contacts of a dedicated starter motor. While the starter-generator dropped out of favour for cars by the 1930s, the concept was still useful for smaller vehicles and was taken up by the German firm SIBA Elektrik which built similar system intended mostly for use on motorcycles, scooters, economy cars (especially those will small-capacity two-stroke engines), and marine engines. These were marketed under the 'Dynastart' name. Since motorcycles usually had small engines and limited electrical equipment, as well as restricted space and weight, the Dynastart was a useful feature. The windings for the starter-generator were usually incorporated into the engine's flywheel, thus not requiring a separate unit at all. The Ford Model T relied on hand cranks until 1919; during the 1920s, electric starters became near-universal on most new cars, making it easier for women and elderly people to drive. It was still common for cars to be supplied with starter handles into the 1960s, and this continued much later for some makes (e.g. Citroën 2CV until end of production in 1990). In many cases, cranks were used for setting timing rather than starting the engine as growing displacements and compression ratios made this impractical. Communist bloc cars such as Ladas often still sported crank-starting as late as the 1980s.
History:
For the first examples of production German turbojet engines later in World War II, Norbert Riedel designed a small two-stroke, opposed-twin gasoline engine to start both the Junkers Jumo 004 and BMW 003 aircraft gas turbines as a form of auxiliary power unit to get the central spindle of each engine design rotating — these were usually installed at the very front of the turbojet, and were themselves started by a pull-rope to get them running during the startup procedure for the jet engines they were fitted to.
History:
Before Chrysler's 1949 innovation of the key-operated combination ignition-starter switch, the starter was often operated by the driver pressing a button mounted on the floor or dashboard. Some vehicles had a pedal in the floor that manually engaged the starter drive pinion with the flywheel ring gear, then completed the electrical circuit to the starter motor once the pedal reached the end of its travel. Ferguson tractors from the 1940s, including the Ferguson TE20, had an extra position on the gear lever that engaged the starter switch, ensuring safety by preventing the tractors from being started in gear.
Electric:
The electric starter motor or cranking motor is the most common type used on gasoline engines and small diesel engines. The modern starter motor is either a permanent-magnet or a series-parallel wound direct current electric motor with a starter solenoid (similar to a relay) mounted on it. When DC power from the starting battery is applied to the solenoid, usually through a key-operated switch (the "ignition switch"), the solenoid engages a lever that pushes out the drive pinion on the starter driveshaft and meshes the pinion with the starter ring gear on the flywheel of the engine.The solenoid also closes high-current contacts for the starter motor, which begins to turn. Once the engine starts, the key-operated switch is opened, a spring in the solenoid assembly pulls the pinion gear away from the ring gear, and the starter motor stops. The starter's pinion is clutched to its drive shaft through an overrunning sprag clutch which permits the pinion to transmit drive in only one direction. In this manner, drive is transmitted through the pinion to the flywheel ring gear, but if the pinion remains engaged (as for example because the operator fails to release the key as soon as the engine starts, or if there is a short and the solenoid remains engaged), the pinion will spin independently of its drive shaft. This prevents the engine driving the starter, for such backdrive would cause the starter to spin so fast as to fly apart.
Electric:
The sprag clutch arrangement would preclude the use of the starter as a generator if employed in the hybrid scheme mentioned above, unless modifications were made. The standard starter motor is typically designed for intermittent use, which would preclude its use as a generator. The starter's electrical components are designed only to operate for typically under 30 seconds before overheating (by too-slow dissipation of heat from ohmic losses), to save weight and cost. Most automobile owner manuals instruct the operator to pause for at least ten seconds after each ten or fifteen seconds of cranking the engine, when trying to start an engine that does not start immediately.
Electric:
This overrunning-clutch pinion arrangement was phased into use beginning in the early 1960s; before that time, a Bendix drive was used. The Bendix system places the starter drive pinion on a helically cut drive shaft. When the starter motor begins turning, the inertia of the drive pinion assembly causes it to ride forward on the helix and thus engage with the ring gear. When the engine starts, backdrive from the ring gear causes the drive pinion to exceed the rotative speed of the starter, at which point the drive pinion is forced back down the helical shaft and thus out of mesh with the ring gear. This has the disadvantage that the gears will disengage if the engine fires briefly but does not continue to run.
Electric:
Folo-Thru drive An intermediate development between the Bendix drive developed in the 1930s and the overrunning-clutch designs introduced in the 1960s was the Bendix Folo-Thru drive. The standard Bendix drive would disengage from the ring gear as soon as the engine fired, even if it did not continue to run. The Folo-Thru drive contains a latching mechanism and a set of flyweights in the body of the drive unit. When the starter motor begins turning and the drive unit is forced forward on the helical shaft by inertia, it is latched into the engaged position. Only once the drive unit is spun at a speed higher than that attained by the starter motor itself (i.e., it is backdriven by the running engine) will the flyweights pull radially outward, releasing the latch and permitting the overdriven drive unit to be spun out of engagement. In this manner, unwanted starter disengagement is avoided before a successful engine start.
Electric:
Gear reduction In 1962, Chrysler introduced a starter incorporating a geartrain between the motor and the drive shaft. The motor shaft included integrally cut gear teeth forming a pinion that meshes with a larger adjacent driven gear to provide a gear reduction ratio of 3.75:1. This permitted the use of a higher-speed, lower-current, lighter and more compact motor assembly while increasing cranking torque. Variants of this starter design were used on most rear- and four-wheel-drive vehicles produced by Chrysler Corporation from 1962 through 1987. It makes a unique, distinct sound when cranking the engine, which led to it being nicknamed the "Highland Park Hummingbird"—a reference to Chrysler's headquarters in Highland Park, Michigan.The Chrysler gear-reduction starter formed the conceptual basis for the gear-reduction starters that now predominate in vehicles on the road. Many Japanese automakers phased in gear reduction starters in the 1970s and 1980s. Light aircraft engines also made extensive use of this kind of starter, because its light weight offered an advantage.
Electric:
Those starters not employing offset gear trains like the Chrysler unit generally employ planetary epicyclic gear trains instead. Direct-drive starters are almost entirely obsolete owing to their larger size, heavier weight and higher current requirements.
Electric:
Movable pole shoe Ford issued a nonstandard starter, a direct-drive "movable pole shoe" design that provided cost reduction rather than electrical or mechanical benefits. This type of starter eliminated the solenoid, replacing it with a movable pole shoe and a separate starter relay. This starter operates as follows: The driver turns the key, activating the starter switch. A small electric current flows through the solenoid actuated starter relay, closing the contacts and sending large battery current to the starter motor. One of the pole shoes, hinged at the front, linked to the starter drive, and spring-loaded away from its normal operating position, is swung into position by the magnetic field created by electricity flowing through its field coil. This moves the starter drive forward to engage the flywheel ring gear, and simultaneously closes a pair of contacts supplying current to the rest of the starter motor winding. Once the engine starts and the driver releases the starter switch, a spring retracts the pole shoe, which pulls the starter drive out of engagement with the ring gear.
Electric:
This starter was used on Ford vehicles from 1973 through 1990, when a gear-reduction unit conceptually similar to the Chrysler unit replaced it.
Electric:
Inertia starter A variant on the electric starter motor is the inertia starter (not to be confused with the Bendix-type starter described above). Here the starter motor does not turn the engine directly. Instead, when energized, the motor turns a heavy flywheel built into its casing (not the main flywheel of the engine). Once the flywheel/motor unit has reached a constant speed the current to the motor is turned off and the drive between the motor and flywheel is disengaged by a freewheel mechanism. The spinning flywheel is then connected to the main engine and its inertia turns it over to start it. These stages are commonly automated by solenoid switches, with the machine operator using a two-position control switch, which is held in one position to spin the motor and then moved to the other to cut the current to the motor and engage the flywheel to the engine.
Electric:
The advantage of the inertia starter is that, because the motor is not driving the engine directly, it can be of much lower power than the standard starter for an engine of the same size. This allows for a motor of much lower weight and smaller size, as well as lighter cables and smaller batteries to power the motor. This made the inertia starter a common choice for aircraft with large radial piston engines. The disadvantage is the increased time required to start the engine - spinning up the flywheel to the required speed can take between 10 and 20 seconds. If the engine does not start by the time the flywheel has lost its inertia then the process must be repeated for the next attempt.
Pneumatic:
Some gas turbine engines and diesel engines, particularly on trucks, use a pneumatic self-starter. In ground vehicles the system consists of a geared turbine, an air compressor and a pressure tank. Compressed air released from the tank is used to spin the turbine, and through a set of reduction gears, engages the ring gear on the flywheel, much like an electric starter. The engine, once running, drives the compressor to recharge the tank.
Pneumatic:
Aircraft with large gas turbine engines are typically started using a large volume of low-pressure compressed air, supplied from a very small engine referred to as an auxiliary power unit, located elsewhere in the aircraft. Alternatively, aircraft gas turbine engines can be rapidly started using a mobile ground-based pneumatic starting engine, referred to as a start cart or air start cart.
Pneumatic:
On larger diesel generators found in large shore installations and especially on ships, a pneumatic starting gear is used. The air motor is normally powered by compressed air at pressures of 10–30 bar. The air motor is made up of a center drum about the size of a soup can with four or more slots cut into it to allow for the vanes to be placed radially on the drum to form chambers around the drum. The drum is offset inside a round casing so that the inlet air for starting is admitted at the area where the drum and vanes form a small chamber compared to the others. The compressed air can only expand by rotating the drum, which allows the small chamber to become larger and puts another one of the cambers in the air inlet. The air motor spins much too fast to be used directly on the flywheel of the engine; instead a large gearing reduction, such as a planetary gear, is used to lower the output speed. A Bendix gear is used to engage the flywheel.
Pneumatic:
Since large trucks typically use air brakes, the system does double duty, supplying compressed air to the brake system. Pneumatic starters have the advantages of delivering high torque, mechanical simplicity and reliability. They eliminate the need for oversized, heavy storage batteries in prime mover electrical systems.
Pneumatic:
Large Diesel generators and almost all Diesel engines used as the prime mover of ships use compressed air acting directly on the cylinder head. This is not ideal for smaller Diesels, as it provides too much cooling on starting. Also, the cylinder head needs to have enough space to support an extra valve for the air start system. The air start system is conceptually very similar to a distributor in a car. There is an air distributor that is geared to the camshaft of the Diesel engine; on the top of the air distributor is a single lobe similar to what is found on a camshaft. Arranged radially around this lobe are roller tip followers for every cylinder. When the lobe of the air distributor hits one of the followers it will send an air signal that acts upon the back of the air start valve located in the cylinder head, causing it to open. Compressed air is provided from a large reservoir that feeds into a header located along the engine. As soon as the air start valve is opened, the compressed air is admitted and the engine will begin turning. It can be used on two-cycle and four-cycle engines and on reversing engines. On large two-stroke engines less than one revolution of the crankshaft is needed for starting.
Hydraulic:
Some Diesel engines from six to 16 cylinders are started by means of a hydraulic motor. Hydraulic starters and the associated systems provide a sparkless, reliable method of engine starting over a wide temperature range. Typically hydraulic starters are found in applications such as remote generators, lifeboat propulsion engines, offshore fire pumping engines, and hydraulic fracturing rigs. The system used to support the hydraulic starter includes valves, pumps, filters, a reservoir, and piston accumulators. The operator can manually recharge the hydraulic system; this cannot readily be done with electric starting systems, so hydraulic starting systems are favored in applications wherein emergency starting is a requirement.
Hydraulic:
With various configurations, Hydraulic starters can be fitted on any engine. Hydraulic starters employ the high efficiency of the axial piston motor concept, which provides high torque at any temperature or environment, and guarantees minimal wear of the engine ring gear and the pinion.
Non-motor:
Spring starter A spring starter uses potential energy stored in a spring wound up with a crank to start an engine without a battery or alternator. Turning the crank moves the pinion into mesh with the engine's ring gear, then winds up the spring. Pulling the release lever then applies the spring tension to the pinion, turning the ring gear to start the engine. The pinion automatically disengages from the flywheel after operation. Provision is also made to allow the engine to be slowly turned over by hand for engine maintenance. This is achieved by operating the trip lever just after the pinion has engaged with the flywheel. Subsequent turning of the winding handle during this operation will not load the starter. Spring starters can be found in engine-generators and hydraulic power packs, and on lifeboat engines, with the most common application being backup starting system on seagoing vessels. Many Briggs & Stratton lawn mowers in the 1960s had hand-cranked spring starters.
Non-motor:
Fuel-starting Some modern gasoline engines with twelve or more cylinders always have at least one or more pistons at the beginning of its power stroke and are able to start by injecting fuel into that cylinder and igniting it. The same procedure can be applied to engines with fewer cylinders, if the engine happens to be stopped at the correct position. This is one way of starting an engine of a car with stop-start system. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Mereology**
Mereology:
In logic, philosophy and related fields, mereology (from Greek μέρος 'part' (root: μερε-, mere-, 'part') and the suffix -logy, 'study, discussion, science') is the study of parts and the wholes they form. Whereas set theory is founded on the membership relation between a set and its elements, mereology emphasizes the meronomic relation between entities, which—from a set-theoretic perspective—is closer to the concept of inclusion between sets.
Mereology:
Mereology has been explored in various ways as applications of predicate logic to formal ontology, in each of which mereology is an important part. Each of these fields provides its own axiomatic definition of mereology. A common element of such axiomatizations is the assumption, shared with inclusion, that the part-whole relation orders its universe, meaning that everything is a part of itself (reflexivity), that a part of a part of a whole is itself a part of that whole (transitivity), and that two distinct entities cannot each be a part of the other (antisymmetry), thus forming a poset. A variant of this axiomatization denies that anything is ever part of itself (irreflexivity) while accepting transitivity, from which antisymmetry follows automatically.
Mereology:
Although mereology is an application of mathematical logic, what could be argued to be a sort of "proto-geometry", it has been wholly developed by logicians, ontologists, linguists, engineers, and computer scientists, especially those working in artificial intelligence. In particular, mereology is also on the basis for a point-free foundation of geometry (see for example the quoted pioneering paper of Alfred Tarski and the review paper by Gerla 1995).
Mereology:
In general systems theory, mereology refers to formal work on system decomposition and parts, wholes and boundaries (by, e.g., Mihajlo D. Mesarovic (1970), Gabriel Kron (1963), or Maurice Jessel (see Bowden (1989, 1998)). A hierarchical version of Gabriel Kron's Network Tearing was published by Keith Bowden (1991), reflecting David Lewis's ideas on gunk. Such ideas appear in theoretical computer science and physics, often in combination with sheaf theory, topos, or category theory. See also the work of Steve Vickers on (parts of) specifications in computer science, Joseph Goguen on physical systems, and Tom Etter (1996, 1998) on link theory and quantum mechanics.
History:
Informal part-whole reasoning was consciously invoked in metaphysics and ontology from Plato (in particular, in the second half of the Parmenides) and Aristotle onwards, and more or less unwittingly in 19th-century mathematics until the triumph of set theory around 1910. Metaphysical ideas of this era that discuss the concepts of parts and wholes include divine simplicity and the classical conception of beauty.
History:
Ivor Grattan-Guinness (2001) sheds much light on part-whole reasoning during the 19th and early 20th centuries, and reviews how Cantor and Peano devised set theory. It appears that the first to reason consciously and at length about parts and wholes was Edmund Husserl, in 1901, in the second volume of Logical Investigations – Third Investigation: "On the Theory of Wholes and Parts" (Husserl 1970 is the English translation). However, the word "mereology" is absent from his writings, and he employed no symbolism even though his doctorate was in mathematics.
History:
Stanisław Leśniewski coined "mereology" in 1927, from the Greek word μέρος (méros, "part"), to refer to a formal theory of part-whole he devised in a series of highly technical papers published between 1916 and 1931, and translated in Leśniewski (1992). Leśniewski's student Alfred Tarski, in his Appendix E to Woodger (1937) and the paper translated as Tarski (1984), greatly simplified Leśniewski's formalism. Other students (and students of students) of Lesniewski elaborated this "Polish mereology" over the course of the 20th century. For a good selection of the literature on Polish mereology, see Srzednicki and Rickey (1984). For a survey of Polish mereology, see Simons (1987). Since 1980 or so, however, research on Polish mereology has been almost entirely historical in nature.
History:
A. N. Whitehead planned a fourth volume of Principia Mathematica, on geometry, but never wrote it. His 1914 correspondence with Bertrand Russell reveals that his intended approach to geometry can be seen, with the benefit of hindsight, as mereological in essence. This work culminated in Whitehead (1916) and the mereological systems of Whitehead (1919, 1920).
History:
In 1930, Henry S. Leonard completed a Harvard Ph.D. dissertation in philosophy, setting out a formal theory of the part-whole relation. This evolved into the "calculus of individuals" of Goodman and Leonard (1940). Goodman revised and elaborated this calculus in the three editions of Goodman (1951). The calculus of individuals is the starting point for the post-1970 revival of mereology among logicians, ontologists, and computer scientists, a revival well-surveyed in Simons (1987), Casati and Varzi (1999), and Cotnoir and Varzi (2021).
Axioms and primitive notions:
Reflexivity: A basic choice in defining a mereological system, is whether to consider things to be parts of themselves. In naive set theory a similar question arises: whether a set is to be considered a "subset" of itself. In both cases, "yes" gives rise to paradoxes analogous to Russell's paradox: Let there be an object O such that every object that is not a proper part of itself is a proper part of O. Is O a proper part of itself? No, because no object is a proper part of itself; and yes, because it meets the specified requirement for inclusion as a proper part of O. In set theory, a set is often termed an improper subset of itself. Given such paradoxes, mereology requires an axiomatic formulation.
Axioms and primitive notions:
A mereological "system" is a first-order theory (with identity) whose universe of discourse consists of wholes and their respective parts, collectively called objects. Mereology is a collection of nested and non-nested axiomatic systems, not unlike the case with modal logic.
The treatment, terminology, and hierarchical organization below follow Casati and Varzi (1999: Ch. 3) closely. For a more recent treatment, correcting certain misconceptions, see Hovda (2008). Lower-case letters denote variables ranging over objects. Following each symbolic axiom or definition is the number of the corresponding formula in Casati and Varzi, written in bold.
Axioms and primitive notions:
A mereological system requires at least one primitive binary relation (dyadic predicate). The most conventional choice for such a relation is parthood (also called "inclusion"), "x is a part of y", written Pxy. Nearly all systems require that parthood partially order the universe. The following defined relations, required for the axioms below, follow immediately from parthood alone: An immediate defined predicate is "x is a proper part of y", written PPxy, which holds (i.e., is satisfied, comes out true) if Pxy is true and Pyx is false. Compared to parthood (which is a partial order), ProperPart is a strict partial order.
Axioms and primitive notions:
PPxy↔(Pxy∧¬Pyx).
3.3 An object lacking proper parts is an atom. The mereological universe consists of all objects we wish to think about, and all of their proper parts:Overlap: x and y overlap, written Oxy, if there exists an object z such that Pzx and Pzy both hold.
Oxy↔∃z[Pzx∧Pzy].
3.1 The parts of z, the "overlap" or "product" of x and y, are precisely those objects that are parts of both x and y.Underlap: x and y underlap, written Uxy, if there exists an object z such that x and y are both parts of z.
Uxy↔∃z[Pxz∧Pyz].
3.2Overlap and Underlap are reflexive, symmetric, and intransitive.
Systems vary in what relations they take as primitive and as defined. For example, in extensional mereologies (defined below), parthood can be defined from Overlap as follows: Pxy↔∀z[Ozx→Ozy].
3.31The axioms are: Parthood partially orders the universe:M1, Reflexive: An object is a part of itself.
Pxx.
P.1 M2, Antisymmetric: If Pxy and Pyx both hold, then x and y are the same object.
(Pxy∧Pyx)→x=y.
P.2 M3, Transitive: If Pxy and Pyz, then Pxz.
(Pxy∧Pyz)→Pxz.
P.3M4, Weak Supplementation: If PPxy holds, there exists a z such that Pzy holds but Ozx does not.
PPxy→∃z[Pzy∧¬Ozx].
P.4M5, Strong Supplementation: If Pyx does not hold, there exists a z such that Pzy holds but Ozx does not.
¬Pyx→∃z[Pzy∧¬Ozx].
P.5M5', Atomistic Supplementation: If Pxy does not hold, then there exists an atom z such that Pzx holds but Ozy does not.
¬Pxy→∃z[Pzx∧¬Ozy∧¬∃v[PPvz]].
P.5' Top: There exists a "universal object", designated W, such that PxW holds for any x.
∃W∀x[PxW].
3.20 Top is a theorem if M8 holds.Bottom: There exists an atomic "null object", designated N, such that PNx holds for any x.
∃N∀x[PNx].
3.22M6, Sum: If Uxy holds, there exists a z, called the "sum" or "fusion" of x and y, such that the objects overlapping of z are just those objects that overlap either x or y.
Uxy→∃z∀v[Ovz↔(Ovx∨Ovy)].
P.6M7, Product: If Oxy holds, there exists a z, called the "product" of x and y, such that the parts of z are just those objects that are parts of both x and y.
Oxy→∃z∀v[Pvz↔(Pvx∧Pvy)].
P.7 If Oxy does not hold, x and y have no parts in common, and the product of x and y is undefined.M8, Unrestricted Fusion: Let φ(x) be a first-order formula in which x is a free variable. Then the fusion of all objects satisfying φ exists.
∃x[ϕ(x)]→∃z∀y[Oyz↔∃x[ϕ(x)∧Oyx]].
Axioms and primitive notions:
P.8 M8 is also called "General Sum Principle", "Unrestricted Mereological Composition", or "Universalism". M8 corresponds to the principle of unrestricted comprehension of naive set theory, which gives rise to Russell's paradox. There is no mereological counterpart to this paradox simply because parthood, unlike set membership, is reflexive.M8', Unique Fusion: The fusions whose existence M8 asserts are also unique. P.8' M9, Atomicity: All objects are either atoms or fusions of atoms.
Various systems:
Simons (1987), Casati and Varzi (1999) and Hovda (2008) describe many mereological systems whose axioms are taken from the above list. We adopt the boldface nomenclature of Casati and Varzi. The best-known such system is the one called classical extensional mereology, hereinafter abbreviated CEM (other abbreviations are explained below). In CEM, P.1 through P.8' hold as axioms or are theorems. M9, Top, and Bottom are optional.
Various systems:
The systems in the table below are partially ordered by inclusion, in the sense that, if all the theorems of system A are also theorems of system B, but the converse is not necessarily true, then B includes A. The resulting Hasse diagram is similar to Fig. 3.2 in Casati and Varzi (1999: 48).
Various systems:
There are two equivalent ways of asserting that the universe is partially ordered: Assume either M1-M3, or that Proper Parthood is transitive and asymmetric, hence a strict partial order. Either axiomatization results in the system M. M2 rules out closed loops formed using Parthood, so that the part relation is well-founded. Sets are well-founded if the axiom of regularity is assumed. The literature contains occasional philosophical and common-sense objections to the transitivity of Parthood.
Various systems:
M4 and M5 are two ways of asserting supplementation, the mereological analog of set complementation, with M5 being stronger because M4 is derivable from M5. M and M4 yield minimal mereology, MM. Reformulated in terms of Proper Part, MM is Simons's (1987) preferred minimal system.
Various systems:
In any system in which M5 or M5' are assumed or can be derived, then it can be proved that two objects having the same proper parts are identical. This property is known as Extensionality, a term borrowed from set theory, for which extensionality is the defining axiom. Mereological systems in which Extensionality holds are termed extensional, a fact denoted by including the letter E in their symbolic names.
Various systems:
M6 asserts that any two underlapping objects have a unique sum; M7 asserts that any two overlapping objects have a unique product. If the universe is finite or if Top is assumed, then the universe is closed under Sum. Universal closure of Product and of supplementation relative to W requires Bottom. W and N are, evidently, the mereological analog of the universal and empty sets, and Sum and Product are, likewise, the analogs of set-theoretical union and intersection. If M6 and M7 are either assumed or derivable, the result is a mereology with closure.
Various systems:
Because Sum and Product are binary operations, M6 and M7 admit the sum and product of only a finite number of objects. The Unrestricted Fusion axiom, M8, enables taking the sum of infinitely many objects. The same holds for Product, when defined. At this point, mereology often invokes set theory, but any recourse to set theory is eliminable by replacing a formula with a quantified variable ranging over a universe of sets by a schematic formula with one free variable. The formula comes out true (is satisfied) whenever the name of an object that would be a member of the set (if it existed) replaces the free variable. Hence any axiom with sets can be replaced by an axiom schema with monadic atomic subformulae. M8 and M8' are schemas of just this sort. The syntax of a first-order theory can describe only a denumerable number of sets; hence, only denumerably many sets may be eliminated in this fashion, but this limitation is not binding for the sort of mathematics contemplated here.
Various systems:
If M8 holds, then W exists for infinite universes. Hence, Top need be assumed only if the universe is infinite and M8 does not hold. Top (postulating W) is not controversial, but Bottom (postulating N) is. Leśniewski rejected Bottom, and most mereological systems follow his example (an exception is the work of Richard Milton Martin). Hence, while the universe is closed under sum, the product of objects that do not overlap is typically undefined. A system with W but not N is isomorphic to: a Boolean algebra lacking a 0; a join semilattice bounded from above by 1. Binary fusion and W interpret join and 1, respectively.Postulating N renders all possible products definable, but also transforms classical extensional mereology into a set-free model of Boolean algebra.
Various systems:
If sets are admitted, M8 asserts the existence of the fusion of all members of any nonempty set. Any mereological system in which M8 holds is called general, and its name includes G. In any general mereology, M6 and M7 are provable. Adding M8 to an extensional mereology results in general extensional mereology, abbreviated GEM; moreover, the extensionality renders the fusion unique. On the converse, however, if the fusion asserted by M8 is assumed unique, so that M8' replaces M8, then—as Tarski (1929) had shown—M3 and M8' suffice to axiomatize GEM, a remarkably economical result. Simons (1987: 38–41) lists a number of GEM theorems.
Various systems:
M2 and a finite universe necessarily imply Atomicity, namely that everything either is an atom or includes atoms among its proper parts. If the universe is infinite, Atomicity requires M9. Adding M9 to any mereological system, X results in the atomistic variant thereof, denoted AX. Atomicity permits economies, for instance, assuming that M5' implies Atomicity and extensionality, and yields an alternative axiomatization of AGEM.
Set theory:
The notion of "subset" in set theory is not entirely the same as the notion of "subpart" in mereology. Stanisław Leśniewski rejected set theory as related to but not the same as nominalism. For a long time, nearly all philosophers and mathematicians avoided mereology, seeing it as tantamount to a rejection of set theory. Goodman too was a nominalist, and his fellow nominalist Richard Milton Martin employed a version of the calculus of individuals throughout his career, starting in 1941.
Set theory:
Much early work on mereology was motivated by a suspicion that set theory was ontologically suspect, and that Occam's razor requires that one minimise the number of posits in one's theory of the world and of mathematics. Mereology replaces talk of "sets" of objects with talk of "sums" of objects, objects being no more than the various things that make up wholes.
Set theory:
Many logicians and philosophers reject these motivations, on such grounds as: They deny that sets are in any way ontologically suspect Occam's razor, when applied to abstract objects like sets, is either a dubious principle or simply false Mereology itself is guilty of proliferating new and ontologically suspect entities such as fusions.For a survey of attempts to found mathematics without using set theory, see Burgess and Rosen (1997).
Set theory:
In the 1970s, thanks in part to Eberle (1970), it gradually came to be understood that one can employ mereology regardless of one's ontological stance regarding sets. This understanding is called the "ontological innocence" of mereology. This innocence stems from mereology being formalizable in either of two equivalent ways: Quantified variables ranging over a universe of sets Schematic predicates with a single free variable.Once it became clear that mereology is not tantamount to a denial of set theory, mereology became largely accepted as a useful tool for formal ontology and metaphysics.
Set theory:
In set theory, singletons are "atoms" that have no (non-empty) proper parts; many consider set theory useless or incoherent (not "well-founded") if sets cannot be built up from unit sets. The calculus of individuals was thought to require that an object either have no proper parts, in which case it is an "atom", or be the mereological sum of atoms. Eberle (1970), however, showed how to construct a calculus of individuals lacking "atoms", i.e., one where every object has a "proper part" (defined below) so that the universe is infinite.
Set theory:
There are analogies between the axioms of mereology and those of standard Zermelo–Fraenkel set theory (ZF), if Parthood is taken as analogous to subset in set theory. On the relation of mereology and ZF, also see Bunt (1985). One of the very few contemporary set theorists to discuss mereology is Potter (2004).
Set theory:
Lewis (1991) went further, showing informally that mereology, augmented by a few ontological assumptions and plural quantification, and some novel reasoning about singletons, yields a system in which a given individual can be both a part and a subset of another individual. Various sorts of set theory can be interpreted in the resulting systems. For example, the axioms of ZFC can be proven given some additional mereological assumptions.
Set theory:
Forrest (2002) revises Lewis's analysis by first formulating a generalization of CEM, called "Heyting mereology", whose sole nonlogical primitive is Proper Part, assumed transitive and antireflexive. There exists a "fictitious" null individual that is a proper part of every individual. Two schemas assert that every lattice join exists (lattices are complete) and that meet distributes over join. On this Heyting mereology, Forrest erects a theory of pseudosets, adequate for all purposes to which sets have been put.
Mathematics:
Husserl never claimed that mathematics could or should be grounded in part-whole rather than set theory. Lesniewski consciously derived his mereology as an alternative to set theory as a foundation of mathematics, but did not work out the details. Goodman and Quine (1947) tried to develop the natural and real numbers using the calculus of individuals, but were mostly unsuccessful; Quine did not reprint that article in his Selected Logic Papers. In a series of chapters in the books he published in the last decade of his life, Richard Milton Martin set out to do what Goodman and Quine had abandoned 30 years prior. A recurring problem with attempts to ground mathematics in mereology is how to build up the theory of relations while abstaining from set-theoretic definitions of the ordered pair. Martin argued that Eberle's (1970) theory of relational individuals solved this problem.
Mathematics:
Topological notions of boundaries and connection can be married to mereology, resulting in mereotopology; see Casati and Varzi (1999: ch. 4,5). Whitehead's 1929 Process and Reality contains a good deal of informal mereotopology.
Natural language:
Bunt (1985), a study of the semantics of natural language, shows how mereology can help understand such phenomena as the mass–count distinction and verb aspect. But Nicolas (2008) argues that a different logical framework, called plural logic, should be used for that purpose.
Natural language:
Also, natural language often employs "part of" in ambiguous ways (Simons 1987 discusses this at length). Hence, it is unclear how, if at all, one can translate certain natural language expressions into mereological predicates. Steering clear of such difficulties may require limiting the interpretation of mereology to mathematics and natural science. Casati and Varzi (1999), for example, limit the scope of mereology to physical objects.
Metaphysics:
In metaphysics there are many troubling questions pertaining to parts and wholes. One question addresses constitution and persistence, another asks about composition.
Metaphysics:
Mereological constitution In metaphysics, there are several puzzles concerning cases of mereological constitution, that is, what makes up a whole. There is still a concern with parts and wholes, but instead of looking at what parts make up a whole, the emphasis is on what a thing is made of, such as its materials, e.g., the bronze in a bronze statue. Below are two of the main puzzles that philosophers use to discuss constitution.
Metaphysics:
Ship of Theseus: Briefly, the puzzle goes something like this. There is a ship called the Ship of Theseus. Over time, the boards start to rot, so we remove the boards and place them in a pile. First question, is the ship made of the new boards the same as the ship that had all the old boards? Second, if we reconstruct a ship using all of the old planks, etc. from the Ship of Theseus, and we also have a ship that was built out of new boards (each added one-by-one over time to replace old decaying boards), which ship is the real Ship of Theseus? Statue and Lump of Clay: Roughly, a sculptor decides to mold a statue out of a lump of clay. At time t1 the sculptor has a lump of clay. After many manipulations at time t2 there is a statue. The question asked is, is the lump of clay and the statue (numerically) identical? If so, how and why?Constitution typically has implications for views on persistence: how does an object persist over time if any of its parts (materials) change or are removed, as is the case with humans who lose cells, change height, hair color, memories, and yet we are said to be the same person today as we were when we were first born. For example, Ted Sider is the same today as he was when he was born—he just changed. But how can this be if many parts of Ted today did not exist when Ted was just born? Is it possible for things, such as organisms to persist? And if so, how? There are several views that attempt to answer this question. Some of the views are as follows (note, there are several other views):(a) Constitution view. This view accepts cohabitation. That is, two objects share exactly the same matter. Here, it follows, that there are no temporal parts.
Metaphysics:
(b) Mereological essentialism, which states that the only objects that exist are quantities of matter, which are things defined by their parts. The object persists if matter is removed (or the form changes); but the object ceases to exist if any matter is destroyed.
(c) Dominant Sorts. This is the view that tracing is determined by which sort is dominant; they reject cohabitation. For example, lump does not equal statue because they're different "sorts".
(d) Nihilism—which makes the claim that no objects exist, except simples, so there is no persistence problem.
(e) 4-dimensionalism or temporal parts (may also go by the names perdurantism or exdurantism), which roughly states that aggregates of temporal parts are intimately related. For example, two roads merging, momentarily and spatially, are still one road, because they share a part.
(f) 3-dimensionalism (may also go by the name endurantism), where the object is wholly present. That is, the persisting object retains numerical identity.
Metaphysics:
Mereological composition One question that is addressed by philosophers is which is more fundamental: parts, wholes, or neither? Another pressing question is called the special composition question (SCQ): For any Xs, when is it the case that there is a Y such that the Xs compose Y? This question has caused philosophers to run in three different directions: nihilism, universal composition (UC), or a moderate view (restricted composition). The first two views are considered extreme since the first denies composition, and the second allows any and all non-spatially overlapping objects to compose another object. The moderate view encompasses several theories that try to make sense of SCQ without saying 'no' to composition or 'yes' to unrestricted composition.
Metaphysics:
Fundamentality There are philosophers who are concerned with the question of fundamentality. That is, which is more ontologically fundamental the parts or their wholes. There are several responses to this question, though one of the default assumptions is that the parts are more fundamental. That is, the whole is grounded in its parts. This is the mainstream view. Another view, explored by Schaffer (2010) is monism, where the parts are grounded in the whole. Schaffer does not just mean that, say, the parts that make up my body are grounded in my body. Rather, Schaffer argues that the whole cosmos is more fundamental and everything else is a part of the cosmos. Then, there is the identity theory which claims that there is no hierarchy or fundamentality to parts and wholes. Instead wholes are just (or equivalent to) their parts. There can also be a two-object view which says that the wholes are not equal to the parts—they are numerically distinct from one another. Each of these theories has benefits and costs associated with them.
Metaphysics:
Special composition question (SCQ) Philosophers want to know when some Xs compose something Y. There are several kinds of responses: One response to this question is called nihilism. Nihilism states that there are no mereological complex objects (read: composite objects); there are only simples. Nihilists do not entirely reject composition because they do think that simples compose themselves, but this is a different point. More formally Nihilists would say: Necessarily, for any non-overlapping Xs, there is an object composed of the Xs if and only if there is only one of the Xs. This theory, though well explored, has its own set of problems. Some of which include, but are not limited to: experiences and common sense, incompatible with atomless gunk, and it is unsupported by space-time physics.
Metaphysics:
Another prominent response is called universal composition (UC). UC says that so long as the Xs do not spatially overlap, the Xs can compose a complex object. Universal compositionalists are also considered those who support unrestricted composition. More formally: Necessarily, for any non-overlapping Xs, there is a Y such that Y is composed of the Xs. For example, someone's left thumb, the top half of another person's right shoe, and a quark in the center of their galaxy can compose a complex object according to universal composition. Likewise, this theory also has some issues, most of them dealing with our experiences that these randomly chosen parts make up a complex whole and there are far too many objects posited in our ontology.
Metaphysics:
A third response (perhaps less explored than the previous two) includes a range of restricted composition views. Though there are several views, they all share a common idea: that there is a restriction on what counts as a complex object: some (but not all) Xs come together to compose a complex Y. Some of these theories include:(a) Contact—the Xs compose a complex Y if and only if the Xs are in contact; (b) Fastenation—the Xs compose a complex Y if and only if the Xs are fastened; (c) Cohesion—the Xs compose a complex Y if and only if the Xs cohere (cannot be pulled apart or moved in relation to each other without breaking); (d) Fusion—the Xs compose a complex Y if and only if the Xs are fused (fusion is when the Xs are joined together such that there is no boundary); (e) Organicism—the Xs compose a complex Y if and only if either the activities of the Xs constitute a life or there is only one of the Xs; and (f) Brutal Composition—"It's just the way things are." There is no true, nontrivial, and finitely long answer.This is not an exhaustive list as many more hypotheses continue to be explored. However, a common problem with these theories is that they are vague. It remains unclear what "fastened" or "life" mean, for example. But there are many other issues within the restricted composition responses—though many of them are subject to which theory is being discussed.
Metaphysics:
A fourth response is called deflationism. Deflationism states that there is variance on how the term "exist" is used, and thus all of the above answers to the SCQ can be correct when indexed to a favorable meaning of "exist." Further, there is no privileged way in which the term "exist" must be used. There is therefore no privileged answer to the SCQ, since there are no privileged conditions for when X composes Y. Instead, the debate is reduced to a mere verbal dispute rather than a genuine ontological debate. In this way, the SCQ is part of a larger debate in general ontological realism and anti-realism. While deflationism successfully avoids the SCQ, it is not devoid of problems. It comes with the cost of ontological anti-realism such that nature has no objective reality at all. For, if there is no privileged way to objectively affirm the existence of objects, nature itself must have no objectivity.
Important surveys:
The books by Simons (1987) and Casati and Varzi (1999) differ in their strengths: Simons (1987) sees mereology primarily as a way of formalizing ontology and metaphysics. His strengths include the connections between mereology and: The work of Stanisław Leśniewski and his descendants Various continental philosophers, especially Edmund Husserl Contemporary English-speaking technical philosophers such as Kit Fine and Roderick Chisholm Recent work on formal ontology and metaphysics, including continuants, occurrents, class nouns, mass nouns, and ontological dependence and integrity Free logic as a background logic Extending mereology with tense logic and modal logic Boolean algebras and lattice theory.
Important surveys:
Casati and Varzi (1999) see mereology primarily as a way of understanding the material world and how humans interact with it. Their strengths include the connections between mereology and: A "proto-geometry" for physical objects Topology and mereotopology, especially boundaries, regions, and holes A formal theory of events Theoretical computer science The writings of Alfred North Whitehead, especially his Process and Reality and work descended therefrom.Simons devotes considerable effort to elucidating historical notations. The notation of Casati and Varzi is often used. Both books include excellent bibliographies. To these works should be added Hovda (2008), which presents the latest state of the art on the axiomatization of mereology.
Sources:
Bowden, Keith, 1991. Hierarchical Tearing: An Efficient Holographic Algorithm for System Decomposition, Int. J. General Systems, Vol. 24(1), pp 23–38.
Bowden, Keith, 1998. Huygens Principle, Physics and Computers. Int. J. General Systems, Vol. 27(1-3), pp. 9–32.
Bunt, Harry, 1985. Mass terms and model-theoretic semantics. Cambridge Univ. Press.
Burgess, John P., and Rosen, Gideon, 1997. A Subject with No Object. Oxford Univ. Press.
Burkhardt, H., and Dufour, C.A., 1991, "Part/Whole I: History" in Burkhardt, H., and Smith, B., eds., Handbook of Metaphysics and Ontology. Muenchen: Philosophia Verlag.
Casati, Roberto, and Varzi, Achille C., 1999. Parts and Places: the structures of spatial representation. MIT Press.
Cotnoir, A. J., and Varzi, Achille C., 2021, Mereology, Oxford University Press.
Eberle, Rolf, 1970. Nominalistic Systems. Kluwer.
Etter, Tom, 1996. Quantum Mechanics as a Branch of Mereology in Toffoli T., et al., PHYSCOMP96, Proceedings of the Fourth Workshop on Physics and Computation, New England Complex Systems Institute.
Etter, Tom, 1998. Process, System, Causality and Quantum Mechanics. SLAC-PUB-7890, Stanford Linear Accelerator Centre.
Forrest, Peter, 2002, "Nonclassical mereology and its application to sets", Notre Dame Journal of Formal Logic 43: 79–94.
Gerla, Giangiacomo, (1995). "Pointless Geometries", in Buekenhout, F., Kantor, W. eds., "Handbook of incidence geometry: buildings and foundations". North-Holland: 1015–31.
Goodman, Nelson, 1977 (1951). The Structure of Appearance. Kluwer.
Goodman, Nelson, and Quine, Willard, 1947, "Steps toward a constructive nominalism", Journal of Symbolic Logic 12: 97-122.
Gruszczynski, R., and Pietruszczak, A., 2008, "Full development of Tarski's geometry of solids", Bulletin of Symbolic Logic 14: 481–540. A system of geometry based on Lesniewski's mereology, with basic properties of mereological structures.
Hovda, Paul, 2008, "What is classical mereology?" Journal of Philosophical Logic 38(1): 55–82.
Husserl, Edmund, 1970. Logical Investigations, Vol. 2. Findlay, J.N., trans. Routledge.
Kron, Gabriel, 1963, Diakoptics: The Piecewise Solution of Large Scale Systems. Macdonald, London.
Lewis, David K., 1991. Parts of Classes. Blackwell.
Leonard, H. S., and Goodman, Nelson, 1940, "The calculus of individuals and its uses", Journal of Symbolic Logic 5: 45–55.
Leśniewski, Stanisław, 1992. Collected Works. Surma, S.J., Srzednicki, J.T., Barnett, D.I., and Rickey, V.F., editors and translators. Kluwer.
Lucas, J. R., 2000. Conceptual Roots of Mathematics. Routledge. Ch. 9.12 and 10 discuss mereology, mereotopology, and the related theories of A.N. Whitehead, all strongly influenced by the unpublished writings of David Bostock.
Mesarovic, M.D., Macko, D., and Takahara, Y., 1970, "Theory of Multilevel, Hierarchical Systems". Academic Press.
Nicolas, David, 2008, "Mass nouns and plural logic", Linguistics and Philosophy 31(2): 211–44.
Pietruszczak, Andrzej, 1996, "Mereological sets of distributive classes", Logic and Logical Philosophy 4: 105–22. Constructs, using mereology, mathematical entities from set theoretical classes.
Pietruszczak, Andrzej, 2005, "Pieces of mereology", Logic and Logical Philosophy 14: 211–34. Basic mathematical properties of Lesniewski's mereology.
Pietruszczak, Andrzej, 2018, Metamerology, Nicolaus Copernicus University Scientific Publishing House.
Potter, Michael, 2004. Set Theory and Its Philosophy. Oxford Univ. Press.
Simons, Peter, 1987 (reprinted 2000). Parts: A Study in Ontology. Oxford Univ. Press.
Srzednicki, J. T. J., and Rickey, V. F., eds., 1984. Lesniewski's Systems: Ontology and Mereology. Kluwer.
Tarski, Alfred, 1984 (1956), "Foundations of the Geometry of Solids" in his Logic, Semantics, Metamathematics: Papers 1923–38. Woodger, J., and Corcoran, J., eds. and trans. Hackett.
Varzi, Achille C., 2007, "Spatial Reasoning and Ontology: Parts, Wholes, and Locations" in Aiello, M. et al., eds., Handbook of Spatial Logics. Springer-Verlag: 945–1038.
Whitehead, A. N., 1916, "La Theorie Relationiste de l'Espace", Revue de Metaphysique et de Morale 23: 423–454. Translated as Hurley, P.J., 1979, "The relational theory of space", Philosophy Research Archives 5: 712–741.
------, 1919. An Enquiry Concerning the Principles of Natural Knowledge. Cambridge Univ. Press. 2nd ed., 1925.
------, 1920. The Concept of Nature. Cambridge Univ. Press. 2004 paperback, Prometheus Books. Being the 1919 Tarner Lectures delivered at Trinity College, Cambridge.
------, 1978 (1929). Process and Reality. Free Press.
Woodger, J. H., 1937. The Axiomatic Method in Biology. Cambridge Univ. Press. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Miniature horse**
Miniature horse:
A miniature horse is a breed or type of horse characterised by its small size. Usually it has been bred to display in miniature the physical characteristics of a full-sized horse, but to be little over 100 cm (40 in) in height, or even less. Although such horses have the appearance of small horses, they are genetically much more similar to pony breeds such as the Shetland.: 486 They have various colors and coat patterns.
Miniature horse:
Miniature horses are present in several countries, including Argentina, Australia, France, Germany, Holland, Ireland, Namibia, the Philippines, the United Kingdom and the United States. In some countries they have the status of a breed; these include the Falabella of Argentina, the Dutch Miniature or Nederlands Mini Paard, the South African Miniature Horse and the American Miniature Horse.
They are commonly kept as companion animals, or for sporting activities such as driving or other competitive horse show events. A few have been trained as guide horses for blind people.
History:
Miniature horses originated in Europe, where there is written and iconographic documentation of them from the late eighteenth century.: 168 In the first half of the twentieth century small horses were bred in England by Lady Estella Mary Hope and her sister Lady Dorothea.: 168 The Falabella was developed in Argentina in the mid-1800s by Patrick Newtall. When Newtall died, the herd and breeding methods were passed to his son-in-law, Juan Falabella. Juan added additional bloodlines including the Welsh Pony, Shetland pony, and small Thoroughbreds. With considerable inbreeding he was able to gain consistently small size within the herd.: 183 The South African Miniature Horse was bred in South Africa from about 1945, when Wynand de Wet of Lindley began selective breeding of Shetland stock. In 2011 there were about 700 of the horses in the country. Morphology is variable: some have an Arab appearance, while others have the conformation of a draft horse. A breed association was established in 1984, and in 1989 the South African Miniature was recognized by the national South African Stud Book and Livestock Improvement Association.
Characteristics:
Miniature horses are generally quite hardy. They often live for longer than is typical for full-sized horses of some breeds; the usual life span is from 25 to 35 years.Their pre-disposition to disease is markedly different from that of full-sized horses. They are only rarely affected by ailments such as laryngeal hemiplegia, osteochondrosis or navicular disease, all of which are common in larger horses,: xii but are much more likely to develop other illnesses rare in large horses, such as hyperlipaemia – which may lead to hepatic lipidosis – or eclampsia.: xii : 6 Dental misalignment and overcrowding are more common than in larger horses: brachygnathism ('parrot mouth') and prognathism ('sow mouth') are often seen;: 53 retention of caps can occur, as can infection of the sinuses associated with tooth eruption.: 55 Poor mastication can contribute to an increased incidence of colic caused by enteroliths, faecoliths, or sand.: 5
Use:
Miniature horses are commonly kept as companion animals. They are often too small for any but the smallest riders to ride, but are well suited to driving, some may participate in other horse show events.: 170 A small number have been trained as guide horses for blind people, particularly for those who consider dogs unclean, as is common in Muslim cultures.: 170 | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Conference on Web and Internet Economics**
Conference on Web and Internet Economics:
Conference on Web and Internet Economics (WINE) (prior to 2013, The Workshop on Internet & Network Economics) is an interdisciplinary workshop devoted to the analysis of algorithmic and economic problems arising in the Internet and the World Wide Web. The submissions are peer reviewed and the proceedings of the conference is published by Springer-Verlag. The conference has been held every year since 2005. Previous sessions include: WINE 2005: Hong Kong, China : Proceedings: Lecture Notes in Computer Science 3828 Springer 2005, ISBN 3-540-30900-4 WINE 2006: Patras, Greece : Proceedings: Lecture Notes in Computer Science 4286 Springer 2006, ISBN 3-540-68138-8 WINE 2007: San Diego, CA, USA : Proceedings: Lecture Notes in Computer Science 4858 Springer 2007, ISBN 978-3-540-77104-3 WINE 2008: Shanghai, China : Proceedings: Lecture Notes in Computer Science 5385 Springer 2008, ISBN 978-3-540-92184-4 WINE 2009: Rome, Italy : Proceedings: Lecture Notes in Computer Science 5929 Springer 2009, ISBN 978-3-642-10840-2 WINE 2010: Stanford, CA, USA : Proceedings: Lecture Notes in Computer Science 6484 Springer 2010, ISBN 978-3-642-17571-8 WINE 2011: Singapore : Proceedings: Lecture Notes in Computer Science 7090 Springer 2011, ISBN 978-3-642-25509-0 WINE 2012: Liverpool, UK : Proceedings: Lecture Notes in Computer Science 7695 Springer 2012, ISBN 978-3-642-35310-9 WINE 2013: Cambridge, MA, USA : Proceedings: Lecture Notes in Computer Science 8289 Springer 2013, ISBN 978-3-642-45046-4 WINE 2014: Beijing, China : Proceedings: Lecture Notes in Computer Science 8877 Springer 2014, ISBN 978-3-319-13128-3 WINE 2015: Amsterdam, The Netherlands : Proceedings: Lecture Notes in Computer Science 9470 Springer 2015, ISBN 978-3-662-48994-9 WINE 2016: Montreal, Canada : Proceedings: Lecture Notes in Computer Science 10123 Springer 2016, ISBN 978-3-662-54109-8 WINE 2017: Bangalore, India : Proceedings: Lecture Notes in Computer Science 10660 Springer 2017, ISBN 978-3-319-71924-5 WINE 2018: Oxford, UK : Proceedings: Lecture Notes in Computer Science 11316 Springer 2018, ISBN 978-3-030-04611-8 WINE 2019: New York, NY, USA : Proceedings: Lecture Notes in Computer Science 11920 Springer 2019, ISBN 978-3-030-35388-9 WINE 2020: Beijing, China : Proceedings: Lecture Notes in Computer Science 12495 Springer 2020, ISBN 978-3-030-64945-6 WINE 2021: Potsdam, Germany : Proceedings: Lecture Notes in Computer Science 13112 Springer 2022, ISBN 978-3-030-94675-3 WINE 2022: Troy, NY, USA : Proceedings: Lecture Notes in Computer Science 13778 Springer 2022, ISBN 978-3-031-22831-5 | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Log reduction**
Log reduction:
Log reduction is a measure of how thoroughly a decontamination process reduces the concentration of a contaminant. It is defined as the common logarithm of the ratio of the levels of contamination before and after the process, so an increment of 1 corresponds to a reduction in concentration by a factor of 10.
Log reduction:
In general, an n-log reduction means that the concentration of remaining contaminants is only 10−n times that of the original. So for example, a 0-log reduction is no reduction at all, while a 1-log reduction corresponds to a reduction of 90 percent from the original concentration, and a 2-log reduction corresponds to a reduction of 99 percent from the original concentration.
Mathematical definition:
Let cb and ca be the numerical values of the concentrations of a given contaminant, respectively before and after treatment, following a defined process.
It is irrelevant in what units these concentrations are given, provided that both use the same units.
Then an and R-log reduction is achieved, where 10 10 10 (cacb) .For the purpose of presentation, the value of R is rounded down to a desired precision, usually to a whole number.
Example Let the concentration of some contaminant be 580 ppm before and 0.725 ppm after treatment. Then 10 0.725 580 10 0.00125 2.903 Rounded down, R is 2, so a 2-log reduction is achieved.
Conversely, an R-log reduction means that a reduction by a factor of 10R has been achieved.
Log reduction and percentage reduction:
Reduction is often expressed as a percentage. The closer it is to 100%, the better.
Letting cb and ca be as before, a reduction by P % is achieved, where 100 ×cb−cacb.
Example Let, as in the earlier example, the concentration of some contaminant be 580 ppm before and 0.725 ppm after treatment. Then 100 580 0.725 580 100 0.99875 99.875.
So this is (better than) a 99% reduction, but not yet quite a 99.9% reduction.
The following table summarizes the most common cases.
In general, if R is a whole number, an R-log reduction corresponds to a percentage reduction with R leading digits "9" in the percentage (provided that it is at least 10%). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Universal Conceptual Cognitive Annotation**
Universal Conceptual Cognitive Annotation:
Universal Conceptual Cognitive Annotation (UCCA) is a semantic approach to grammatical representation. It is a cross-linguistically applicable semantic representation scheme, and has demonstrated support for rapid annotation. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Voronezhselmash**
Voronezhselmash:
Voronezhselmash (Russian: Воронежсельмаш) produces equipment for post-harvest handling, drying and storing grain, including grain elevators and separators. Construction of grain elevators for turnkey grain storage. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Three mountain problem**
Three mountain problem:
The Three Mountains Task was a task developed by Jean Piaget, a developmental psychologist from Switzerland. Piaget came up with a theory for developmental psychology based on cognitive development. Cognitive development, according to his theory, took place in four stages. These four stages were classified as the sensorimotor, preoperational, concrete operational and formal operational stages. The Three Mountain Problem was devised by Piaget to test whether a child's thinking was egocentric, which was also a helpful indicator of whether the child was in the preoperational stage or the concrete operational stage of cognitive development.
Methods:
Piaget's aim in the Three Mountain Problem was to investigate egocentrism in children's thinking. The original setup for the task was: The child who is seated at a table where a model of three mountains is presented in front. The mountains were of different sizes, and they had different identifiers (one mountain had snow; one had a red cross on top; one had a hut on top). The child was allowed to do a 360 surveillance of the model. Upon having a good look at the model, a doll is placed at different vantage points relative to the child, and the child is shown 10 photographs. The child is to select which of the 10 photographs best reflects the doll's view. Children of different ages were tested using this task to determine the age at which children begin to 'decenter,' or take the perspective of others.
Findings:
The findings showed that at age 4, children would choose the photograph that best reflected with their own view. At age 6, an awareness of perspective different from their own could be seen. Then, by ages 7–8, children can clearly acknowledge more than one point of view and consistently select the correct photograph.
During Preoperational Stage:
A distinction can be made between children who are in the preoperational stage of cognitive development and the concrete operational stage. The prototypical child in the preoperational stage will fail the Three Mountain Problem task. The child will choose the photograph that best represents their own viewpoint, not that of the doll's.
During Preoperational Stage:
What is implied is that the child's selection is based on egocentric thinking. Egocentric thinking is looking at the world from the child's point of view solely, thus "an egocentric child assumes that other people see, hear, and feel exactly the same as the child does.” This is consistent with the results for the preoperational age range as they selected photographs paralleling their own view.
During Preoperational Stage:
On a similar note, these results help Piaget hone in on what age children show the capacity to decenter their thoughts, otherwise seen in a deviation away from egocentric thinking. Preoperational children have not achieved this yet; their thinking is centered, which is defined as a propensity to focus on one salient aspect or one dimension of a problem while simultaneously neglecting other potentially relevant aspects.
During Concrete Operational Stage:
The concept of centration is observed predominantly in children in the preoperational stage of cognitive development. Conversely, children in the concrete operational stage demonstrate decentration - an ability to recognize alternate point of views and a straying away from egocentric thinking. Piaget concluded that, by age 7, children were able to decenter their thoughts and acknowledge perspectives different than their own. This was evidenced by the consistent and correct selection of photographs by seven- and eight-year-olds in the 1956 study.
During Concrete Operational Stage:
An example of a correct answer would be if the child and the doll were situated on the complete opposite sides of the mountain model with a tree on the child's side and a large mountain in the middle acting as a visual barrier. A preoperational child would claim that the doll could see the tree, whereas the concrete operational child would select a photograph without the tree since the mountain is large enough to block the tree from the doll's view. A concrete operational child would pass the Three Mountain Problem task.
Follow Up Studies:
There has been some criticism that the Three Mountain Problem was too difficult for the children to understand, compounded with the additional requirement of matching their answer to a photograph. Martin Hughes conducted a study in 1975 called the Policeman Doll Study. Two intersecting walls were used to create different quadrants, and “policeman" dolls were moved in various locations. The children were asked to hide another doll, a “boy” doll, away from both policemen's views. The results showed that among the sample of children ranging from ages 3.5-5, 90% gave correct answers. When the stakes were raised and additional walls and policeman dolls were added, 90% of four-year-olds were still able to pass the task. Hughes claimed that because this task made more sense to the child (with a primer session with one police doll to guarantee this), children were able to exhibit a loss of egocentric thinking as early as four years of age.
Variations of the Three Mountain Problem:
Common criticism of the Three Mountain Problem is about the complexity of the task. In 1975, another researcher by the name of Helen Borke replicated the task using a farm area with landmarks such as a lake, animals, people, trees, and a building. A character from Sesame Street, Grover, was put in a car, and he was driven around the area. When he stopped to "take a look at the scenery," children were asked what the landscape looked like from Grover's perspective. The results showed that children as young as three-years-old were able to perform well, and they showed evidence of perspective-taking, the ability to understand a situation from an alternate point of view. Hence, evaluation of Piaget's Three Mountain Problem has shown that using objects more familiar to the child and making the task less complex will produce different results than the original study. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**LE (text editor)**
LE (text editor):
LE is a text editor which appears something like the Norton Editor, but has many additional features: Rectangle select/copy/paste (block type is switchable) Search/replace with regular expressions Filtering block contents through an external program Linear multilevel undo/redo Customizable menus Color syntax highlighting (using regular expressions in an external file) Handles UTF-8 characters, based on locale settings Customizable keymaps for different terminal types (associating either literal strings or terminfo capability names) Hexadecimal editing mode Editing of mmap'd files or devices in replace mode Frame drawing mode (first seen in Lexicon) File selection box (inspired by Turbo C) Built-in postfix calculator.It uses ncurses for display, mouse and part of the keyboard handling. The application has a built-in table of key assignments for xterm, rxvt and some less familiar terminal types.
History:
According to the HISTORY file in its sources, Alexander V. Lukyanov started writing it in 1993 while using a BESTA machine. Over the next four years, he rewrote it into C++, and published it in 1997 under the GNU General Public License. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**IAd Producer**
IAd Producer:
iAd Producer was introduced by Apple Inc. in 2010 as a new authoring tool for designing and developing interactive iAd using HTML5, CSS3, and JavaScript standards for distribution through its iAd network within iOS apps. iAd and by extension iAd Producer were both discontinued in June 2016.
IAd Producer:
iAd Producer offers a single-window interface for viewing and editing layouts and settings of iAd rich media ad projects. It gives each access to project templates with pre-built structure including banners, splash pages, and menus. An extensive library of pre-built interactive elements – carousels, galleries, maps, videos and more – which are available for use using simple drag and drop. Sophisticated object animation tools are available to use through the GUI using a timeline that is very easy to manage as it only shows events on the timeline and not the dead time in between them. The app also doubles as an IDE, by offering advanced JavaScript code editing and debugging with full syntax highlighting and code completion. One can even write their own plugin to provide additional functionality for use in their projects. For testing purposes, iAd Producer leverages Mac OS X, Safari, the iOS simulator and hardware iOS devices.
Version History:
Since version 2.0, Apple has maintained an archive of iAd Producer release notes. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Graphene lens**
Graphene lens:
A graphene lens is an optical refraction device. Graphene's unique 2-D honeycomb contributes to its unique optical properties.
Graphene:
The honeycomb structure allows electrons to behave as massless quasiparticles known as Dirac fermions. Graphene's optical conductivity properties are thus unobstructed by any material parameters, as represented by equation 1, where e is the electron charge, h is Planck's constant and e2/h represents the universal conductance.
Graphene:
uni =(π)e22h (Equation 1) This behavior is the result of an undoped graphene material at zero temperature (figure 1a). In contrast to traditional semiconductors or metals (figure 1b); graphene's band gap is nearly nonexistent because the conducting and valence bands make contact (Figure 1a). However, the band gap is tunable via doping and electrical gating, changing optical properties. As a result of its tunable conductivity, graphene is suitable for various optical applications.
Applications:
Photodetectors Electrical gating and doping allows for adjustment of graphene's optical absorptivity. The application of electric fields transverse to staggered graphene bilayers generates a shift in Fermi energy and an artificial, non-zero band gap (equation 2 figure 1).
Applications:
δD=Dt−Db (Equation 2)where Dt = top electrical displacement field Db = bottom electrical displacement fieldVarying δD above or below zero (δD=0 denotes non-gated, neutral bilayers) allows electrons to pass through the bilayer without altering the gating-induced band gap. As shown in Figure 2, varying the average displacement field, ▁D, alters the bilayer's absorption spectra. The optical tunability resulting from gating and electrostatic doping (also known as charge plasma doping) lends to the application of graphene as an ultra-broadband photodetector in lenses.
Applications:
Chang-Hua et al. implemented graphene in an infrared photodetector by sandwiching an insulating barrier of Ta2O5 between two graphene sheets. The graphene layers became electrically isolated and exhibited an average Fermi difference of 0.12eV when a current was passed through the bottom layer (Figure 3). When the photodetector is exposed to light, excited hot electrons transitioned from the top graphene layer to the bottom, a process promoted by the structural asymmetry of the insulating Ta2O5 barrier. As a consequence of the hot electron transition, the top layer accumulates positive charges and induces a photogating effect on the lower graphene layer, which is measured as a change in current correlating with photon detection. Utilizing graphene both as a channel for charge transport and light absorption, the photodetectors ably detects the visible to mid-infrared spectrum. Nanometers thin and functional at room temperature, graphene photodetectors show promise in lens applications.
Applications:
Fresnel zone plates Fresnel zone plates are devices that focus light on a fixed point in space. These devices concentrate light reflected off a lens onto a singular point (Figure 4). Composed of a series of discs centered about an origin, Fresnel zone plates are manufactured using laser pulses, which embed voids into areflective lens.
Applications:
Despite its weak reflectance (R = 0.25π2 α 2 at T = 1.3 × 10-4 K), graphene has utility as a lens for Fresnel zone plates. Graphene lenses effectively concentrate light of ʎ = 850 nm onto a single point 120 um away from the Fresnel zone plate (figure 5). Further investigation illustrates that the reflected intensity increases linearly with the number of graphene layers within the lens (Figure 6).
Applications:
Transparent conductors Optoelectronic components such as light-emitting diode (LED) displays, solar cells, and touchscreens require highly transparent materials with low sheet resistance, Rs. For a thin film, the sheet resistance is given by Equation 3: Rs=tσ (Equation 3)where t is the film thickness and σ is the DC conductivity.
Applications:
A material with tunable thickness t and conductivity σ is suitable for optoelectronic applications if Rs is reasonably small. Graphene is such a material; the number of graphene layers that comprise the film can tune t and the inherent tunability of graphene's optical properties via doping or grating can tune sigma. Figure 7 shows graphene's potential relative to other known transparent conductors.
Applications:
The need for alternative transparent conductors is well documented. Semiconductor based transparent conductors such as doped indium oxides, zinc oxides, or tin oxides suffer from practical downfalls including rigorous processing requirements, prohibitive cost, sensitivity toward Ph, and brittle consistency. However, graphene does not suffer from these shortfalls. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**ACS Energy Letters**
ACS Energy Letters:
ACS Energy Letters is a monthly peer-reviewed scientific journal published by the American Chemical Society. It was established in 2016 and the editor-in-chief is Prashant V. Kamat (University of Notre Dame). It covers research on all aspects of energy and aims for rapid publication.
Abstracting and indexing:
The journal is abstracted and indexed in: Ei Compendex Current Contents/Engineering, Computing & Technology Current Contents/Physical, Chemical & Earth Sciences Inspec Science Citation Index Expanded ScopusAccording to the Journal Citation Reports, the journal has a 2022 impact factor of 22.
Article types:
The journal publishes the following article types: letters, energy express, reviews, perspectives, viewpoints, energy focus, and editorials. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Equivariant stable homotopy theory**
Equivariant stable homotopy theory:
In mathematics, more specifically in topology, the equivariant stable homotopy theory is a subfield of equivariant topology that studies a spectrum with group action instead of a space with group action, as in stable homotopy theory. The field has become more active recently because of its connection to algebraic K-theory. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Polysome**
Polysome:
A polyribosome (or polysome or ergosome) is a group of ribosomes bound to an mRNA molecule like “beads” on a “thread”. It consists of a complex of an mRNA molecule and two or more ribosomes that act to translate mRNA instructions into polypeptides. Originally coined "ergosomes" in 1963, they were further characterized by Jonathan Warner, Paul M. Knopf, and Alex Rich.
Polysome:
Polysomes are formed during the elongation phase when ribosomes and elongation factors synthesize the encoded polypeptide. Multiple ribosomes move along the coding region of mRNA, creating a polysome. The ability of multiple ribosomes to function on an mRNA molecule explains the limited abundance of mRNA in the cell. Polyribosome structure differs between prokaryotic polysomes, eukaryotic polysomes, and membrane bound polysomes. Polysome activity can be used to measure the level of gene expression through a technique called polysomal profiling.
Structure:
Electron microscopy technologies such as staining, metal shadowing, and ultra-thin cell sections were the original methods to determine polysome structure. The development of cryo-electron microscopy techniques has allowed for increased resolution of the image, leading to a more precise method to determine structure. Different structural configurations of polyribosomes could reflect a variety in translation of mRNAs. An investigation of the ratio of polyribosomal shape elucidated that a high number of circular and zigzag polysomes were found after several rounds of translation. A longer period of translation caused the formation of densely packed 3-D helical polysomes. Different cells produce different structures of polysomes.
Structure:
Prokaryotic Bacterial polysomes have been found to form double-row structures. In this conformation, the ribosomes are contacting each other through smaller subunits. These double row structures generally have a “sinusoidal” (zigzag) or 3-D helical path. In the “sinusoidal” path, there are two types of contact between the small subunits- “top-to-top” or “top-to-bottom”. In the 3-D helical path, only “top-to-top” contact is observed.Polysomes are present in archaea, but not much is known about the structure.
Structure:
Eukaryotic In cells in situ (in cell) studies have shown that eukaryotic polysomes exhibit linear configurations. Densely packed 3-D helices and planar double-row polysomes were found with variable packing including “top-to-top” contacts similar to prokaryotic polysomes. Eukaryotic 3-D polyribosomes are similar to prokaryotic 3-D polyribosomes in that they are “densely packed left-handed helices with four ribosomes per turn”. This dense packing can determine their function as regulators of translation, with 3-D polyribosomes being found in sarcoma cells using fluorescence microscopy.
Structure:
Cell free Atomic force microscopy used in in vitro studies have shown that circular eukaryotic polysomes can be formed by free polyadenylated mRNA in the presence of initiation factor eIF4E bound to the 5’ cap and PABP bound to the 3’-poly(A) tail. However, this interaction between cap and poly(A)-tail mediated by a protein complex is not a unique way of circularizing polysomal mRNA. It has been found that topologically circular polyribosomes can be successfully formed in the translational system with mRNA with no cap and no poly(A) tail as well as a capped mRNA without a 3’-poly(A) tail.
Structure:
Membrane-bound Polyribosomes bound to membranes are restricted by a 2 dimensional space given by the membrane surface. The restriction of inter-ribosomal contacts causes a round-shape configuration that arranges ribosomes along the mRNA so that the entry and exit sites form a smooth pathway. Each ribosome is turned relative to the previous one, resembling a planar spiral.
Profiling:
Polysomal Profiling is a technique that uses cycloheximide to arrest translation and a sucrose gradient to separate the resulting cell extract by centrifugation. Ribosome-associated mRNAs migrate faster than free mRNAs and polysome associated mRNAs migrate faster than ribosome associated mRNAs. Several peaks corresponding to mRNA are revealed by the measurement of total protein across the gradient. The corresponding mRNA is associated with increasing numbers of ribosomes as polysomes. The presence of mRNA across the gradient reveals the translation of the mRNA. Polysomal profiling is optimally applied to cultured cells and tissues to track the translational status of an identified mRNA as well as measure ribosome density. This technique has been used to compare the translational status of mRNAs in different cell types.
Profiling:
For example, polysomal profiling was used in a study to investigate the effect of vesicular stomatitis virus (VSV) in mammalian cells. The data from polysomal profiling showed that host mRNAs are outcompeted by viral mRNAs for polysomes, therefore decreasing the translation of host mRNA and increasing the translation of viral mRNA. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Roberts linkage**
Roberts linkage:
A Roberts linkage is a four-bar linkage which converts a rotational motion to approximate straight-line motion.The linkage was developed by Richard Roberts.The Roberts linkage can be classified as: Watt-type linkage Grashof rocker-rocker Symmetrical four-bar linkage | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Design closure**
Design closure:
Design Closure is a part of the digital electronic design automation workflow by which an integrated circuit (i.e. VLSI) design is modified from its initial description to meet a growing list of design constraints and objectives.
Every step in the IC design (such as static timing analysis, placement, routing, and so on) is already complex and often forms its own field of study. This article, however, looks at the overall design closure process, which takes a chip from its initial design state to the final form in which all of its design constraints are met.
Introduction:
Every chip starts off as someone’s idea of a good thing: "If we can make a part that performs function X, we will all be rich!" Once the concept is established, someone from marketing says "To make this chip profitably, it must cost $C and run at frequency F." Someone from manufacturing says "To meet this chip’s targets, it must have a yield of Y%." Someone from packaging says “It must fit in the P package and dissipate no more than W watts.” Eventually, the team generates an extensive list of all the constraints and objectives they must meet to manufacture a product that can be sold profitably. The management then forms a design team, which consists of chip architects, logic designers, functional verification engineers, physical designers, and timing engineers, and assigns them to create a chip to the specifications.
Introduction:
Constraints vs Objectives The distinction between constraints and objectives is straightforward: a constraint is a design target that must be met for the design to be successful. For example, a chip may be required to run at a specific frequency so it can interface with other components in a system. In contrast, an objective is a design target where more (or less) is better. For example, yield is generally an objective, which is maximized to lower manufacturing cost. For the purposes of design closure, the distinction between constraints and objectives is not important; this article uses the words interchangeably.
Evolution of the Design Closure Flow:
Designing a chip used to be a much simpler task. In the early days of VLSI, a chip consisted of a few thousand logic circuits that performed a simple function at speeds of a few MHz. Design closure was simple: if all of the necessary circuits and wires "fit", the chip would perform the desired function.
Evolution of the Design Closure Flow:
Modern design closure has grown orders of magnitude more complex. Modern logic chips can have tens to hundreds of millions of logic elements switching at speeds of several GHz. This improvement has been driven by Moore’s law of scaling of technology, and has introduced many new design considerations. As a result, a modern VLSI designer must consider the performance of a chip against a list of dozens of design constraints and objectives including performance, power, signal integrity, reliability, and yield. In response to this growing list of constraints, the design closure flow has evolved from a simple linear list of tasks to a very complex, highly iterative flow such as the following simplified ASIC design flow:
Reference ASIC Design Flow:
Concept phase: Functional objectives and architecture of a chip are developed.
Logic design: Architecture is implemented in a register transfer level (RTL) language, then simulated to verify that it performs the desired functions. This includes functional verification.
Floorplanning: The RTL of the chip is assigned to gross regions of the chip, input/output (I/O) pins are assigned and large objects (arrays, cores, etc.) are placed.
Logic synthesis: The RTL is mapped into a gate-level netlist in the target technology of the chip.
Design for Testability: The test structures like scan chains are inserted.
Placement: The gates in the netlist are assigned to nonoverlapping locations on the chip.
Logic/placement refinement: Iterative logical and placement transformations to close performance and power constraints.
Clock insertion: Balanced buffered clock trees are introduced into the design.
Routing: The wires that connect the gates in the netlist are added.
Postwiring optimization: Remaining performance, noise, and yield violations are removed.
Design for manufacturability: The design is modified, where possible, to make it as easy as possible to produce.
Signoff checks: Since errors are expensive, time consuming and hard to spot, extensive error checking is the rule, making sure the mapping to logic was done correctly, and checking that the manufacturing rules were followed faithfully.
Tapeout and mask generation: the design data is turned into photomasks in mask data preparation.
Evolution of design constraints:
The purpose of the flow is to take a design from concept phase to working chip. The complexity of the flow is a direct result of the addition and evolution of the list of design closure constraints. To understand this evolution it is important to understand the life cycle of a design constraint. In general, design constraints influence the design flow via the following five-stage evolution: Early warnings: Before chip issues begin occurring, academics and industry visionaries make dire predictions about the future impact of some new technology effect.
Evolution of design constraints:
Hardware problems: Sporadic hardware failures start showing up in the field due to the new effect. Postmanufacturing redesign and hardware re-spins are required to get the chip to function.
Trial and error: Constraints on the effect are formulated and used to drive postdesign checking. Violations of the constraint are fixed manually.
Find and repair: Large number of violations of the constraint drives the creation of automatic postdesign analysis and repair flows.
Evolution of design constraints:
Predict and prevent: Constraint checking moves earlier in the flow using predictive estimations of the effect. These drive optimizations to prevent violations of the constraint.A good example of this evolution can be found in the signal integrity constraint. In the mid-1990s (180 nm node), industry visionaries were describing the impending dangers of coupling noise long before chips were failing. By the mid-late 1990s, noise problems were cropping up in advanced microprocessor designs.
Evolution of design constraints:
By 2000, automated noise analysis tools were available and were used to guide manual fix-up. The total number of noise problems identified by the analysis tools identified by the flow quickly became too many to correct manually. In response, CAD companies developed the noise avoidance flows that are currently in use in the industry.
Evolution of design constraints:
At any point in time, the constraints in the design flow are at different stages of their life cycle. At the time of this writing, for example, performance optimization is the most mature and is well into the fifth phase with the widespread use of timing-driven design flows. Power- and defect-oriented yield optimization is well into the fourth phase; power supply integrity, a type of noise constraint, is in the third phase; circuit-limited yield optimization is in the second phase, etc. A list of the first-phase impending constraint crises can always be found in the International Technology Roadmap for Semiconductors (ITRS) 15-year-outlook technology roadmaps.
Evolution of design constraints:
As a constraint matures in the design flow, it tends to work its way from the end of the flow to the beginning.
As it does this, it also tends to increase in complexity and in the degree that it contends with other constraints.
Constraints tend to move up in the flow due to one of the basic paradoxes of design: accuracy vs.
influence. Specifically, the earlier in a design flow a constraint is addressed, the more flexibility there is to address the constraint. Ironically, the earlier one is in a design flow, the more difficult it is to predict compliance.
Evolution of design constraints:
For example, an architectural decision to pipeline a logic function can have a far greater impact on total chip performance than any amount of postrouting fix-up. At the same time, accurately predicting the performance impact of such a change before the chip logic is synthesized, let alone placed or routed, is very difficult. This paradox has shaped the evolution of the design closure flow in several ways. First, it requires that the design flow is no longer composed of a linear set of discrete steps. In the early stages of VLSI it was sufficient to break the design into discrete stages, i.e., first do logic synthesis, then do placement, then do routing. As the number and complexity of design closure constraints has increased, the linear design flow has broken down. In the past if there were too many timing constraint violations left after routing, it was necessary to loop back, modify the tool settings slightly, and reexecute the previous placement steps. If the constraints were still not met, it was necessary to reach further back in the flow and modify the chip logic and repeat the synthesis and placement steps. This type of looping is both time consuming and unable to guarantee convergence i.e., it is possible to loop back in the flow to correct one constraint violation only to find that the correction induced another unrelated violation. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Shell builtin**
Shell builtin:
In computing, a shell builtin is a command or a function, called from a shell, that is executed directly in the shell itself, instead of an external executable program which the shell would load and execute.Shell builtins work significantly faster than external programs, because there is no program loading overhead. However, their code is inherently present in the shell, and thus modifying or updating them requires modifications to the shell. Therefore, shell builtins are usually used for simple, almost trivial, functions, such as text output.
Shell builtin:
Because of the nature of some operating systems, some functions of the systems must necessarily be implemented as shell builtins. The most notable example is the cd command, which changes the working directory of the shell. Since each executable program runs in a separate process, and working directories are specific to each process, loading cd as an external program would not affect the working directory of the shell that loaded it. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Measurement Science and Technology**
Measurement Science and Technology:
Measurement Science and Technology is a monthly peer-reviewed scientific journal, published by IOP Publishing, covering the areas of measurement, instrumentation, and sensor technology in the sciences. The editor-in-chief is Andrew Yacoot (National Physical Laboratory).
History:
The journal was established in 1923 as the Journal of Scientific Instruments. The first issue was introduced by J. J. Thomson, then president of the Institute of Physics, who stated that no publication existed at that time in the English language specially devoted to scientific instruments. The idea for the journal was promoted by Richard Glazebrook, the first president, then director, of the National Physical Laboratory, where the journal was initially edited. The need for interdisciplinarity was recognised even then, with the desire to co-opt biologists, engineers, chemists, and instrument makers, "as well as physicists", on the scientific advisory committee. The Institute of Physics merged with the Physical Society of London in 1960. By this time the Proceedings of the Physical Society had grown in size and the quality of the applied journals, British Journal of Applied Physics and Journal of Scientific Instruments, had been improved. In 1968 these journals were merged to form part of the Journal of Physics series of journals, A to E, the fifth journal in the series being Journal of Physics E: Scientific Instruments. In 1990 the journal was renamed as Measurement Science and Technology to reflect the shift away from many scientists making their own instruments. Since 2003 the journal archive containing all articles published since 1874 are available online.
Abstracting and indexing:
The journal is abstracted and indexed in: According to the website, the journal has a 2020 impact factor of 2.046. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Ansys HFSS**
Ansys HFSS:
Ansys HFSS (high-frequency structure simulator), is a commercial finite element method solver for electromagnetic (EM) structures from Ansys that offers multiple state-of-the-art solver technologies. Each solver in ANSYS HFSS is an automated solution processor for which the user dictates the geometry, properties of the material and the required range of solution frequencies.Engineers use Ansys HFSS primarily to design and simulate high-speed, high-frequency electronics in radar systems, communication systems, satellites, ADAS, microchips, printed circuit boards, IoT products, and other digital devices and RF devices. The solver has also been used to simulate the electromagnetic behavior of objects such as automobiles and aircraft. ANSYS HFSS allows system and circuit designers to simulate EM issues such as losses due to attenuation, coupling, radiation and reflection.The benefits of simulating a circuit's high frequency behavior with high accuracy on a computer reduces the final testing and verification effort of the system as well as mitigating the necessity of building costly multiple prototypes, saving both time and money in product development.HFSS captures and simulates objects in 3D, accounting for materials composition and shapes/geometries of each object. HFSS is one of several commercial tools used for antenna design, and the design of complex radio frequency electronic circuit elements including filters, transmission lines, and packaging.
History:
HFSS was originally developed by Professor Zoltan Cendes, Ph.D., and his students at Carnegie Mellon University. It was the first general purpose software product to solve arbitrary 3D EM field problems, including EM energy distribution and S parameters in complex structures.
History:
In 1984, Dr. Cendes founded Ansoft Corporation to design and develop high performance EDA software. He served as its chairman and chief technology officer until 2008, when Ansys acquired Ansoft.Ansoft originally sold HFSS as a stand-alone product under an agreement with Hewlett-Packard. It was also bundled into Ansoft products.In 1997 Hewlett-Packard acquired Optimization Systems Associates Inc. (OSA), a company John Bandler founded in 1983. HP's acquisition was driven by HP's need for an optimization capability for HFSS. After various business relationships over the period 1996–2006, HP (which became Agilent EEsof EDA division) and Ansoft went their separate ways:Over time, Ansys HFSS introduced a number of new technologies in computational EM simulation, including automatic adaptive mesh generation, tangential vector finite elements, transfinite elements, and reduced-order modeling. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Cowboy bowline**
Cowboy bowline:
The cowboy bowline or left-hand bowline, is a variation of the bowline loop knot. The cowboy bowline has the working end go around the standing part on the side closer to the loop and results with the working end outside the loop. In contrast, a regular bowline has the working end finishing inside the loop. (The "rabbit" goes around the "tree" in the opposite direction from normal.) The Ashley Book of Knots states that it is "distinctly inferior" to the standard bowline because of its similarity to the left-hand sheet bend. Various tests of the different versions' strengths show little difference; conjecture about either knot's vulnerability to some failure remain pretty much only that – conjectures. However, the left-hand bowline is much more stable under ring loading, as it then acts effectively as a proper Lapp bend, while the simple bowline acts as the inferior version of the Lapp bend, which tends to slip.
Cowboy bowline:
As for the tail of a regular bowline finishing "inside the loop [eye]", that is more a formal-image state than one of an actual knot, as the draw of the standing part will pull the tail around so that it actually points away from the eye, but there are various ways the knot can be dressed to affect this aspect.
Cowboy bowline:
Some hearsay suggest the Dutch Navy uses (or used) this variant of the bowline because they consider it superior since the working end is not so easily pushed back by accident. However, there is no documentation to confirm this claim, and some Dutch knot tyers outright deny it.
Another piece of unverified lore is that it is called a winter bowline because the exposed working end on the outside would blow in the wind and prevent it from freezing to the loop on ships in the north Atlantic during winter. (This suggests that the standard bowline would be the summer bowline.)
Security:
There is a rule of thumb which states that the loose end should be as long as 12 times the circumference of the cord for the sake of safety.Cowboy bowline is tested to be more resistant to cross loading (ring loading, transverse loading) than regular bowline. Like the right-hand bowline, it can spontaneously loosen under cyclic loading, and is not recommended for life-critical applications. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Congenital afibrinogenemia**
Congenital afibrinogenemia:
Congenital afibrinogenemia is a rare, genetically inherited blood fibrinogen disorder in which the blood does not clot normally due to the lack of fibrinogen, a blood protein necessary for coagulation. This disorder is autosomal recessive, meaning that two unaffected parents can have a child with the disorder. The lack of fibrinogen expresses itself with excessive and, at times, uncontrollable bleeding.
Signs and symptoms:
As this is a disorder that is present in an individual from birth, there are no warning signs to look for. The first symptom usually seen is hemorrhage from the umbilical cord that is difficult to stop.Other symptoms include: Nasal and oral mucosa bleeds Gastrointestinal bleeding Excessive/spontaneous bleeding or bruising from minor injury Prolonged menstruation in women Spontaneous abortion during pregnancy CNS hemorrhagingSpontaneous bleeding of the mouth, nose, and gastrointestinal tract are common. Since blood clots can not be formed, minor injuries tend to lead to excessive bleeding or bruising. The biggest concern for individuals with afibrinogenemia is CNS hemorrhage, bleeding in the brain, which can be fatal.
Signs and symptoms:
Many of these symptoms are chronic, and will continue to occur for the entirety of the affected individual's life.
Causes:
A missense or nonsense mutation to the genes that code for the fibrinogen protein are affected. Usually the mutation leads to an early stop in the production of the protein. Due to the problem being genetically based, there is no way to prevent the disease. Individuals can get genetic testing done to see if they are a carrier of the trait, and if so may choose to complete genetic counseling to better understand the disorder and help manage family planning. Parents can choose to do prenatal genetic testing for the disorder to determine if their child will have the disease. The only risk factor is if both parents of a child carry the recessive allele linked to the disorder.
Mechanism:
Individuals with the disorder have a mutation to their fibrinogen gene that prevents the formation of the protein. In normal conditions, fibrinogen is converted to fibrin when it is cleaved by the enzyme thrombin in the blood. The newly formed fibrin forms a fiber-rich network that helps trap red blood cells to start the coagulation process and form a clot. Since there is no fibrinogen present during afibrinogenemia, fibrin can not be formed to aid in this process.
Diagnosis:
Diagnostic tests When a problem of fibrinogen is suspected, the following tests can be ordered: PT PTT Fibrinogen level in blood (total and clottable) Reptilase time Thrombin timeBlood fibrinogen levels of less than 0.1 g/L and prolonged bleeding test times are indicators of an individual having afibrinogenemia.A suspicion of congenital afibrinogenemia may appear on a platelet aggregation function test.
Treatment:
The most common treatment is fibrinogen replacement therapy, including cryoprecipitate, blood plasma, and fibrinogen concentration transfusions to help with bleeding episodes or prior to surgery. Although some thrombotic complications have been reported following replacement therapy, transfusions of fibrinogen concentrate are widely considered to be the most beneficial. Fibrinogen concentrate is pure, contains a known quantity of fibrinogen, is virally inactivated, and is transfused in small amounts. There is a lower chance of allergic reaction in response to transfusion. Alternatively, cryoprecipitate contains other coagulation factors, factors VIII, XIII, and von Willebrand factor. There are no known cures or forms of holistic care to date. Most complications arise from the symptoms of the disorder.
Prognosis:
As there is not much data out on the life expectancy of an individual with afibrinogenemia, it is difficult to determine the average lifespan. However, the leading cause of death thus far is linked to CNS hemorrhage and postoperative bleeding.
History:
It was first described in 1920 by German doctors, Fritz Rabe and Eugene Salomon, studying a bleeding disorder presenting itself in a child from birth. This disorder may also be simply called afibrinogenemia or familial afibrinogenemia. About 1 in 1 million individuals are diagnosed with the disease; typically at birth. Both males and females seem to be affected equally, but it has a higher occurrence in regions where consanguinity is prevalent.
Research:
Currently research is based in pharmacological treatments. A case from 2015 was seen in which congenital afibrinogenemia was resolved in a patient after receiving a liver transplant. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Air base**
Air base:
An airbase (stylised air base in American English), sometimes referred to as a military airbase, military airfield, military airport, air station, naval air station, air force station, or air force base, is an aerodrome or airport used as a military base by a military force for the operation of military aircraft.
Airbase facilities:
An airbase typically has some facilities similar to a civilian airport; for example, air traffic control and firefighting. Some military aerodromes have passenger facilities; for example, RAF Brize Norton in England has a terminal used by passengers for the Royal Air Force's passenger transport flights. A number of military airbases may also have a civil enclave for commercial passenger flights, e.g. Beijing Nanyuan Airport (China), Chandigarh Airport (India), Ibaraki Airport (Japan), Burlington International Airport (USA), Sheikh Ul-Alam International Airport Srinagar (India), Taipei Songshan Airport (Taiwan), Eindhoven airport (The Netherlands). Likewise, the opposite also occurs; large civilian airports may contain a smaller military airbase within their environs, such as Royal Brunei Air Force Base, Rimba (located within Brunei International Airport).
Airbase facilities:
Some airbases have dispersed aircraft parking, revetments, hardened aircraft shelters, or even underground hangars, to protect aircraft from enemy attack. Combat aircraft require secure protected storage of aircraft ordnance and munitions. Other facilities may also include technical buildings for servicing and support of survival equipment (including flying helmets and personal liquid oxygen), flight simulator for synthetic training, servicing facilities for all aircraft systems (airframes, propulsion, avionics, weapons systems) and associated ground support systems (including mechanical transport). All military airbases will have buildings for military administration (station headquarters, squadron briefing and operations), and larger bases will also include medical and dental facilities for military personnel (and sometimes their dependents), along with dining (mess, informally known as the 'cook house'), accommodation (single living accommodation for junior ranks, Sergeants' and Officers' Mess for senior non-commissioned officers and commissioned officers), recreational facilities (club house for socialising), shopping facilities (NAAFI shops, base exchange, commissary), and sports facilities (gymnasium, swimming pool, sports pitches). An airbase may be defended by anti-aircraft weapons and force protection troops.
Dispersal airbase:
A dispersal (or dispersed) airbase is an airfield that is used for the purpose of dispersing air units in the event of conflict, so to minimise the vulnerability of aircraft and its supporting units whilst on the ground. Dispersal airbases are not necessarily ordinarily operational in peace time, and may only be activated when needed. Airfields used as dispersal bases can either be auxiliary military airfields, civilian airports, or highway strips. Examples of uses of dispersal bases are the Swedish Bas 60 and Bas 90 systems, the British V-Bomber dispersal bases, and NATO's Dispersed Operating Bases in France.
Road airbase:
Road airbases are highways constructed to double as auxiliary airbases in the event of war. Nations known to utilise this strategy are India, Sweden, Finland, Germany (formerly), Singapore, Switzerland, South Korea, Turkey, Poland, Pakistan, and the Czech Republic. In the case of Finnish road airbases, the space needed for landing aircraft is reduced by means of an arrestor wire, similar to that used on some aircraft carriers (Finnish Air Force uses F-18s, which can land on aircraft carriers).
Aircraft carrier:
An aircraft carrier is a type of naval ship which serves as a seaborne airbase, the development of which has greatly enhanced the capabilities of modern air forces and naval aviators. In many countries, they are now a key part of the military, allowing for their military aircraft to be staged much nearer the area of conflict. Aircraft carriers were vital to the United States during World War II, and to the United Kingdom in the 1982 Falklands War. They retain modern roles as well as "several acres of sovereign territory a nation can move about at will", which allows greater flexibility in diplomacy as well as military affairs. Aircraft carriers may also used in disaster relief. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Vault (architecture)**
Vault (architecture):
In architecture, a vault (French voûte, from Italian volta) is a self-supporting arched form, usually of stone or brick, serving to cover a space with a ceiling or roof. As in building an arch, a temporary support is needed while rings of voussoirs are constructed and the rings placed in position. Until the topmost voussoir, the keystone, is positioned, the vault is not self-supporting. Where timber is easily obtained, this temporary support is provided by centering consisting of a framed truss with a semicircular or segmental head, which supports the voussoirs until the ring of the whole arch is completed.
Vault types:
Corbelled vaults, also called false vaults, with horizontally joined layers of stone have been documented since prehistoric times; in the 14th century BC from Mycenae. They were built regionally until modern times.
The real vault construction with radially joined stones was already known to the Egyptians and Assyrians and was introduced into the building practice of the West by the Etruscans. The Romans in particular developed vault construction further and built barrel, cross and dome vaults. Some outstanding examples have survived in Rome, e.g. the Pantheon and the Basilica of Maxentius.
Vault types:
Brick vaults have been used in Egypt since the early 3rd millennium BC. widely used and from the end of the 8th century B.C. Keystone vaults were built. However, monumental temple buildings of the pharaonic culture in the Nile Valley did not use vaults, since even the huge portals with widths of more than 7 meters were spanned with cut stone beams.
Vault types:
Dome Amongst the earliest known examples of any form of vaulting is to be found in the neolithic village of Khirokitia on Cyprus. Dating from c. 6000 BCE, the circular buildings supported beehive shaped corbel domed vaults of unfired mud-bricks and also represent the first evidence for settlements with an upper floor. Similar beehive tombs, called tholoi, exist in Crete and Northern Iraq. Their construction differs from that at Khirokitia in that most appear partially buried and make provision for a dromos entry.
Vault types:
The inclusion of domes, however, represents a wider sense of the word vault. The distinction between the two is that a vault is essentially an arch which is extruded into the third dimension, whereas a dome is an arch revolved around its vertical axis.
Vault types:
Pitched brick barrel vault Pitched-brick vaults are named for their construction, the bricks are installed vertically (not radially) and are leaning (pitched) at an angle: This allows their construction to be completed without the use of centering. Examples have been found in archaeological excavations in Mesopotamia dating to the 2nd and 3rd millennium BCE, which were set in gypsum mortar.
Vault types:
Barrel vault A barrel vault is the simplest form of a vault and resembles a barrel or tunnel cut lengthwise in half. The effect is that of a structure composed of continuous semicircular or pointed sections.The earliest known examples of barrel vaults were built by the Sumerians, possibly under the ziggurat at Nippur in Babylonia, which was built of fired bricks cemented with clay mortar.The earliest barrel vaults in ancient Egypt are thought to be those in the granaries built by the 19th dynasty Pharaoh Ramesses II, the ruins of which are behind the Ramesseum, at Thebes. The span was 12 feet (3.7 m) and the lower part of the arch was built in horizontal courses, up to about one-third of the height, and the rings above were inclined back at a slight angle, so that the bricks of each ring, laid flatwise, adhered till the ring was completed, no centering of any kind being required; the vault thus formed was elliptic in section, arising from the method of its construction. A similar system of construction was employed for the vault over the great hall at Ctesiphon, where the material employed was fired bricks or tiles of great dimensions, cemented with mortar; but the span was close upon 83 feet (25 m), and the thickness of the vault was nearly 5 feet (1.5 m) at the top, there being four rings of brickwork.Assyrian palaces used pitched-brick vaults, made with sun-dried mudbricks, for gates, subterranean graves and drains. During the reign of king Sennacherib they were used to construct aqueducts, such as those at Jerwan. In the provincial city Dūr-Katlimmu they were used to created vaulted platforms. The tradition of their erection, however, would seem to have been handed down to their successors in Mesopotamia, viz. to the Sassanians, who in their palaces in Sarvestan and Firouzabad built domes of similar form to those shown in the Nimrud sculptures, the chief difference being that, constructed in rubble stone and cemented with mortar, they still exist, though probably abandoned on the Islamic invasion in the 7th century.
Vault types:
Groin vaults A groin vault is formed by the intersection of two or more barrel vaults, resulting in the formation of angles or groins along the lines of transition between the webs. In these bays the longer transverse arches are semi-circular, as are the shorter longitudinal arches. The curvatures of these bounding arches were apparently used as the basis for the web centrings, which was created in the form of two intersecting tunnels as though each web was an arch projected horizontally in three dimensions.The earliest example is thought to be over a small hall at Pergamum, in Asia Minor, but its first employment over halls of great dimensions is due to the Romans. When two semicircular barrel vaults of the same diameter cross one another their intersection (a true ellipse) is known as a groin vault, down which the thrust of the vault is carried to the cross walls; if a series of two or more barrel vaults intersect one another, the weight is carried on to the piers at their intersection and the thrust is transmitted to the outer cross walls; thus in the Roman reservoir at Baiae, known as the Piscina Mirabilis, a series of five aisles with semicircular barrel vaults are intersected by twelve cross aisles, the vaults being carried on 48 piers and thick external walls. The width of these aisles being only about 13 feet (4.0 m) there was no great difficulty in the construction of these vaults, but in the Roman Baths of Caracalla the tepidarium had a span of 80 feet (24 m), more than twice that of an English cathedral, so that its construction both from the statical and economical point of view was of the greatest importance.
Vault types:
The researches of M. Choisy (L'Art de bâtir chez les Romains), based on a minute examination of those portions of the vaults which still remain in situ, have shown that, on a comparatively slight centering, consisting of trusses placed about 10 feet (3.0 m) apart and covered with planks laid from truss to truss, were laid – to begin with – two layers of the Roman brick (measuring nearly 2 feet (0.61 m) square and 2 in. thick); on these and on the trusses transverse rings of brick were built with longitudinal ties at intervals; on the brick layers and embedding the rings and cross ties concrete was thrown in horizontal layers, the haunches being filled in solid, and the surface sloped on either side and covered over with a tile roof of low pitch laid direct on the concrete. The rings relieved the centering from the weight imposed, and the two layers of bricks carried the concrete till it had set.
Vault types:
As the walls carrying these vaults were also built in concrete with occasional bond courses of brick, the whole structure was homogeneous. One of the important ingredients of the mortar was a volcanic deposit found near Rome, known as pozzolana, which, when the concrete had set, not only made the concrete as solid as the rock itself, but to a certain extent neutralized the thrust of the vaults, which formed shells equivalent to that of a metal lid; the Romans, however, do not seem to have recognized the value of this pozzolana mixture, for they otherwise provided amply for the counteracting of any thrust which might exist by the erection of cross walls and buttresses. In the tepidaria of the Thermae and in the basilica of Constantine, in order to bring the thrust well within the walls, the main barrel vault of the hall was brought forward on each side and rested on detached columns, which constituted the principal architectural decoration. In cases where the cross vaults intersecting were not of the same span as those of the main vault, the arches were either stilted so that their soffits might be of the same height, or they formed smaller intersections in the lower part of the vault; in both of these cases, however, the intersections or groins were twisted, for which it was very difficult to form a centering, and, moreover, they were of disagreeable effect: though every attempt was made to mask this in the decoration of the vault by panels and reliefs modelled in stucco.
Vault types:
Rib vault A rib vault is one in which all of the groins are covered by ribs or diagonal ribs in the form of segmental arches. Their curvatures are defined by the bounding arches. Whilst the transverse arches retain the same semi-circular profile as their groin-vaulted counterparts, the longitudinal arches are pointed with both arcs having their centres on the impost line. This allows the latter to correspond more closely to the curvatures of the diagonal ribs, producing a straight tunnel running from east to west.Reference has been made to the rib vault in Roman work, where the intersecting barrel vaults were not of the same diameter. Their construction must at all times have been somewhat difficult, but where the barrel vaulting was carried round over the choir aisle and was intersected (as in St Bartholomew-the-Great in Smithfield, London) by semicones instead of cylinders, it became worse and the groins more complicated. This would seem to have led to a change of system and to the introduction of a new feature, which completely revolutionized the construction of the vault. Hitherto the intersecting features were geometrical surfaces, of which the diagonal groins were the intersections, elliptical in form, generally weak in construction and often twisting. The medieval builder reversed the process, and set up the diagonal ribs first, which were utilized as permanent centres, and on these he carried his vault or web, which henceforward took its shape from the ribs. Instead of the elliptical curve which was given by the intersection of two semicircular barrel vaults, or cylinders, he employed the semicircular arch for the diagonal ribs; this, however, raised the centre of the square bay vaulted above the level of the transverse arches and of the wall ribs, and thus gave the appearance of a dome to the vault, such as may be seen in the nave of Sant'Ambrogio, Florence. To meet this, at first the transverse and wall ribs were stilted, or the upper part of their arches was raised, as in the Abbaye-aux-Hommes at Caen, and the Abbey of Lessay, in Normandy. The problem was ultimately solved by the introduction of the pointed arch for the transverse and wall ribs – the pointed arch had long been known and employed, on account of its much greater strength and of the less thrust it exerted on the walls. When employed for the ribs of a vault, however narrow the span might be, by adopting a pointed arch, its summit could be made to range in height with the diagonal rib; and, moreover, when utilized for the ribs of the annular vault, as in the aisle round the apsidal termination of the choir, it was not necessary that the half ribs on the outer side should be in the same plane as those of the inner side; for when the opposite ribs met in the centre of the annular vault, the thrust was equally transmitted from one to the other, and being already a broken arch the change of its direction was not noticeable.
Vault types:
The first introduction of the pointed arch rib took place at Cefalù Cathedral and pre-dated the abbey of Saint-Denis. Whilst the pointed rib-arch is often seen as an identifier for Gothic architecture, Cefalù is a Romanesque cathedral whose masons experimented with the possibility of Gothic rib-arches before it was widely adopted by western church architecture. Besides Cefalù Cathedral, the introduction of the pointed arch rib would seem to have taken place in the choir aisles of the abbey of Saint-Denis, near Paris, built by the abbot Suger in 1135. It was in the church at Vezelay (1140) that it was extended to the square bay of the porch. As has been pointed out, the aisles had already in the early Christian churches been covered over with groined vaults, the only advance made in the later developments being the introduction of transverse ribs' dividing the bays into square compartments. In the 12th century the first attempts were made to vault over the naves, which were twice the width of the aisles, so it became necessary to include two bays of the aisles to form one rectangular bay in the nave (although this is often mistaken as square). It followed that every alternate pier served no purpose, so far as the support of the nave vault was concerned, and this would seem to have suggested an alternative to provide a supplementary rib across the church and between the transverse ribs. This resulted in what is known as a sexpartite, or six-celled vault, of which one of the earliest examples is found in the Abbaye-aux-Hommes at Caen. This church, built by William the Conqueror, was originally constructed to carry a timber roof only, but nearly a century later the upper part of the nave walls were partly rebuilt, in order that it might be covered with a vault. The immense size, however, of the square vault over the nave necessitated some additional support, so that an intermediate rib was thrown across the church, dividing the square compartment into six cells, and called the sexpartite vault The intermediate rib, however, had the disadvantage of partially obscuring one side of the clerestory windows, and it threw unequal weights on the alternate piers, so that in the cathedral of Soissons (1205) a quadripartite or four-celled vault was introduced, the width of each bay being half the span of the nave, and corresponding therefore with the aisle piers. To this there are some exceptions, in Sant' Ambrogio, Milan, and San Michele, Pavia (the original vault), and in the cathedrals of Speyer, Mainz and Worms, where the quadripartite vaults are nearly square, the intermediate piers of the aisles being of much smaller dimensions. In England sexpartite vaults exist at Canterbury (1175) (set out by William of Sens), Rochester (1200), Lincoln (1215), Durham (east transept), and St. Faith's chapel, Westminster Abbey.In the earlier stage of rib vaulting, the arched ribs consisted of independent or separate voussoirs down to the springing; the difficulty, however, of working the ribs separately led to two other important changes: (1) the lower part of the transverse diagonal and wall ribs were all worked out of one stone; and (2) the lower horizontal, constituting what is known as the tas-de-charge or solid springer. The tas-de-charge, or solid springer, had two advantages: (1) it enabled the stone courses to run straight through the wall, so as to bond the whole together much better; and (2) it lessened the span of the vault, which then required a centering of smaller dimensions. As soon as the ribs were completed, the web or stone shell of the vault was laid on them. In some English work each course of stone was of uniform height from one side to the other; but, as the diagonal rib was longer than either the transverse or wall rib, the courses dipped towards the former, and at the apex of the vault were cut to fit one another. In the early English Gothic period, in consequence of the great span of the vault and the very slight rise or curvature of the web, it was thought better to simplify the construction of the web by introducing intermediate ribs between the wall rib and the diagonal rib and between the diagonal and the transverse ribs; and in order to meet the thrust of these intermediate ribs a ridge rib was required, and the prolongation of this rib to the wall rib hid the junction of the web at the summit, which was not always very sightly, and constituted the ridge rib. In France, on the other hand, the web courses were always laid horizontally, and they are therefore of unequal height, increasing towards the diagonal rib. Each course also was given a slight rise in the centre, so as to increase its strength; this enabled the French masons to dispense with the intermediate rib, which was not introduced by them till the 15th century, and then more as a decorative than a constructive feature, as the domical form given to the French web rendered unnecessary the ridge rib, which, with some few exceptions, exists only in England. In both English and French vaulting centering was rarely required for the building of the web, a template (Fr. cerce) being employed to support the stones of each ring until it was complete. In Italy, Germany and Spain the French method of building the web was adopted, with horizontal courses and a domical form. Sometimes, in the case of comparatively narrow compartments, and more especially in clerestories, the wall rib was stilted, and this caused a peculiar twisting of the web, where the springing of the wall rib is at K: to these twisted surfaces the term ploughshare vaulting is given.
Vault types:
One of the earliest examples of the introduction of the intermediate rib is found in the nave of Lincoln Cathedral, and there the ridge rib is not carried to the wall rib. It was soon found, however, that the construction of the web was much facilitated by additional ribs, and consequently there was a tendency to increase their number, so that in the nave of Exeter Cathedral three intermediate ribs were provided between the wall rib and the diagonal rib. In order to mask the junction of the various ribs, their intersections were ornamented with richly carved bosses, and this practice increased on the introduction of another short rib, known as the lierne, a term in France given to the ridge rib. Lierne ribs are short ribs crossing between the main ribs, and were employed chiefly as decorative features, as, for instance, in the Liebfrauenkirche (1482) of Mühlacker, Germany. One of the best examples of Lierne ribs exists in the vault of the oriel window of Crosby Hall, London. The tendency to increase the number of ribs led to singular results in some cases, as in the choir of Gloucester Cathedral, where the ordinary diagonal ribs become mere ornamental mouldings on the surface of an intersected pointed barrel vault, and again in the cloisters, where the introduction of the fan vault, forming a concave-sided conoid, returned to the principles of the Roman geometrical vault. This is further shown in the construction of these fan vaults, for although in the earliest examples each of the ribs above the tas-de-charge was an independent feature, eventually it was found easier to carve them and the web out of the solid stone, so that the rib and web were purely decorative and had no constructional or independent functions.
Vault types:
Fan vault This form of vaulting is found in English late Gothic in which the vault is constructed as a single surface of dressed stones, with the resulting conoid forming an ornamental network of blind tracery.The fan vault would seem to have owed its origin to the employment of centerings of one curve for all the ribs, instead of having separate centerings for the transverse, diagonal wall and intermediate ribs; it was facilitated also by the introduction of the four-centred arch, because the lower portion of the arch formed part of the fan, or conoid, and the upper part could be extended at pleasure with a greater radius across the vault. These ribs were often cut from the same stones as the webs, with the entire vault being treated as a single jointed surface covered in interlocking tracery.The earliest example is perhaps the east walk of the cloister at Gloucester, with its surface consisting of intricately decorated panels of stonework forming conical structures that rise from the springers of the vault. In later examples, as in King's College Chapel, Cambridge, on account of the great dimensions of the vault, it was found necessary to introduce transverse ribs, which were required to give greater strength. Similar transverse ribs are found in Henry VII's chapel and in the Divinity School at Oxford, where a new development presented itself. One of the defects of the fan vault at Gloucester is the appearance it gives of being half sunk in the wall; to remedy this, in the two buildings just quoted, the complete conoid is detached and treated as a pendant.
Byzantine vaults and domes:
The vault of the Basilica of Maxentius, completed by Constantine, was the last great work carried out in Rome before its fall, and two centuries pass before the next important development is found in the Church of the Holy Wisdom (Hagia Sophia) at Constantinople. It is probable that the realization of the great advance in the science of vaulting shown in this church owed something to the eastern tradition of dome vaulting seen in the Assyrian domes, which are known to us only by the representations in the bas-relief from Nimrud, because in the great water cisterns in Istanbul, known as the Basilica Cistern and Bin bir direk (cistern with a thousand and one columns), we find the intersecting groin vaults of the Romans already replaced by small cupolas or domes. These domes, however, are of small dimensions when compared with that projected and carried out by Justinian in the Hagia Sophia. Previous to this the greatest dome was that of the Pantheon at Rome, but this was carried on an immense wall 20 feet (6.1 m) thick, and with the exception of small niches or recesses in the thickness of the wall could not be extended, so that Justinian apparently instructed his architect to provide an immense hemicycle or apse at the eastern end, a similar apse at the western end, and great arches on either side, the walls under which would be pierced with windows. Unlike the Pantheon dome, the upper portions of which are made of concrete, Byzantine domes were made of brick, which were lighter and thinner, but more vulnerable to the forces exerted onto them.
Byzantine vaults and domes:
The diagram shows the outlines of the solution of the problem. If a hemispherical dome is cut by four vertical planes, the intersection gives four semicircular arches; if cut in addition by a horizontal plane tangent to the top of these arches, it describes a circle; that portion of the sphere which is below this circle and between the arches, forming a spherical spandrel, is the pendentive, and its radius is equal to the diagonal of the square on which the four arches rest. Having obtained a circle for the base of the dome, it is not necessary that the upper portion of the dome should spring from the same level as the arches, or that its domical surface should be a continuation of that of the pendentive. The first and second dome of the Hagia Sophia apparently fell down, so that Justinian determined to raise it, possibly to give greater lightness to the structure, but mainly in order to obtain increased light for the interior of the church. This was effected by piercing it with forty windows – the effect of which, as the light streaming through these windows, gave the dome the appearance of being suspended in the air. The pendentive which carried the dome rested on four great arches, the thrust of those crossing the church being counteracted by immense buttresses which traversed the aisles, and the other two partly by smaller arches in the apse, the thrust being carried to the outer walls, and to a certain extent by the side walls which were built under the arches. From the description given by Procopius we gather that the centering employed for the great arches consisted of a wall erected to support them during their erection. The construction of the pendentives is not known, but it is surmised that to the top of the pendentives they were built in horizontal courses of brick, projecting one over the other, the projecting angles being cut off afterwards and covered with stucco in which the mosaics were embedded; this was the method employed in the erection of the Périgordian domes, to which we shall return; these, however, were of less diameter than those of the Hagia Sophia, being only about 40 to 60 feet (18 m) instead of 107 feet (33 m) The apotheosis of Byzantine architecture, in fact, was reached in Hagia Sophia, for although it formed the model on which all subsequent Byzantine churches were based, so far as their plan was concerned, no domes approaching the former in dimensions were even attempted. The principal difference in some later examples is that which took place in the form of the pendentive on which the dome was carried. Instead of the spherical spandril of Hagia Sophia, large niches were formed in the angles, as in the Mosque of Damascus, which was built by Byzantine workmen for the Al-Walid I in CE 705; these gave an octagonal base on which the hemispherical dome rested; or again, as in the Sassanian palaces of Sarvestan and Firouzabad of the 4th and 5th century, when a series of concentric arch rings, projecting one in front of the other, were built, giving also an octagonal base; each of these pendentives is known as a squinch.
Byzantine vaults and domes:
There is one other remarkable vault, also built by Justinian, in the Church of the Saints Sergius and Bacchus in Constantinople. The central area of this church was octagonal on plan, and the dome is divided into sixteen compartments; of these eight consist of broad flat bands rising from the centre of each of the walls, and the alternate eight are concave cells over the angles of the octagon, which externally and internally give to the roof the appearance of an umbrella.
Romanesque:
Although the dome constitutes the principal characteristic of the Byzantine church, throughout Asia Minor are numerous examples in which the naves are vaulted with the semicircular barrel vault, and this is the type of vault found throughout the south of France in the 11th and 12th centuries, the only change being the occasional substitution of the pointed barrel vault, adopted not only on account of its exerting a less thrust, but because, as pointed out by Fergusson (vol. ii. p. 46), the roofing tiles were laid directly on the vault and a less amount of filling in at the top was required.The continuous thrust of the barrel vault in these cases was met either by semicircular or pointed barrel vaults on the aisles, which had only half the span of the nave; of this there is an interesting example in the Chapel of Saint John in the Tower of London – and sometimes by half-barrel vaults. The great thickness of the walls, however, required in such constructions would seem to have led to another solution of the problem of roofing over churches with incombustible material, viz. that which is found throughout Périgord and La Charente, where a series of domes carried on pendentives covered over the nave, the chief peculiarities of these domes being the fact that the arches carrying them form part of the pendentives, which are all built in horizontal courses.The intersecting and groined vault of the Romans was employed in the early Christian churches in Rome, but only over the aisles, which were comparatively of small span, but in these there was a tendency to raise the centres of these vaults, which became slightly domical; in all these cases centering was employed.
Gothic Revival and the Renaissance:
One good example of the fan vault is that over the staircase leading to the hall of Christ Church, Oxford, where the complete conoid is displayed in its centre carried on a central column. This vault, not built until 1640, is an example of traditional workmanship, probably in Oxford transmitted in consequence of the late vaulting of the entrance gateways to the colleges. Fan vaulting is peculiar to England, the only example approaching it in France being the pendant of the Lady-chapel at Caudebec-en-Caux, in Normandy.In France, Germany, and Spain the multiplication of ribs in the 15th century led to decorative vaults of various kinds, but with some singular modifications. Thus, in Germany, recognizing that the rib was no longer a necessary constructive feature, they cut it off abruptly, leaving a stump only; in France, on the other hand, they gave still more importance to the rib, by making it of greater depth, piercing it with tracery and hanging pendants from it, and the web became a horizontal stone paving laid on the top of these decorated vertical webs. This is the characteristic of the great Renaissance work in France and Spain; but it soon gave way to Italian influence, when the construction of vaults reverted to the geometrical surfaces of the Romans, without, however, always that economy in centering to which they had attached so much importance, and more especially in small structures. In large vaults, where it constituted an important expense, the chief boast of some of the most eminent architects has been that centering was dispensed with, as in the case of the dome of the Santa Maria del Fiore in Florence, built by Filippo Brunelleschi, and Ferguson cites as an example the great dome of the church at Mousta in Malta, erected in the first half of the 19th century, which was built entirely without centering of any kind.
Vaulting and faux-vaulting in the Renaissance and after:
It is important to note that whereas Roman vaults, like that of the Pantheon, and Byzantine vaults, like that at Hagia Sophia, were not protected from above (i.e. the vault from the inside was the same that one saw from the outside), the European architects of the Middle Ages protected their vaults with wooden roofs. In other words, one will not see a Gothic vault from the outside. The reasons for this development are hypothetical, but the fact that the roofed basilica form preceded the era when vaults begin to be made is certainly to be taken into consideration. In other words, the traditional image of a roof took precedence over the vault.
Vaulting and faux-vaulting in the Renaissance and after:
The separation between interior and exterior – and between structure and image – was to be developed very purposefully in the Renaissance and beyond, especially once the dome became reinstated in the Western tradition as a key element in church design. Michelangelo's dome for St. Peter's Basilica in Rome, as redesigned between 1585 and 1590 by Giacomo della Porta, for example, consists of two domes of which, however, only the inner is structural. Baltasar Neumann, in his baroque churches, perfected light-weight plaster vaults supported by wooden frames. These vaults, which exerted no lateral pressures, were perfectly suited for elaborate ceiling frescoes. In St Paul's Cathedral in London there is a highly complex system of vaults and faux-vaults. The dome that one sees from the outside is not a vault, but a relatively light-weight wooden-framed structure resting on an invisible – and for its age highly original – catenary vault of brick, below which is another dome, (the dome that one sees from the inside), but of plaster supported by a wood frame. From the inside, one can easily assume that one is looking at the same vault that one sees from the outside.
India:
There are two distinctive "other ribbed vaults" (called "Karbandi" in Persian) in India which form no part of the development of European vaults, but have some unusual features; one carries the central dome of the Jumma Musjid at Bijapur (A.D. 1559), and the other is Gol Gumbaz, the tomb of Muhammad Adil Shah II (1626–1660) in the same town. The vault of the latter was constructed over a hall 135 feet (41 m) square, to carry a hemispherical dome. The ribs, instead of being carried across the angles only, thus giving an octagonal base for the dome, are carried across to the further pier of the octagon and consequently intersect one another, reducing the central opening to 97 feet (30 m) in diameter, and, by the weight of the masonry they carry, serving as counterpoise to the thrust of the dome, which is set back so as to leave a passage about 12 feet (3.7 m) wide round the interior. The internal diameter of the dome is 124 feet (38 m), its height 175 feet (53 m) and the ribs struck from four centres have their springing 57 feet (17 m) from the floor of the hall. The Jumma Musjid dome was of smaller dimensions, on a square of 70 feet (21 m) with a diameter of 57 feet (17 m), and was carried on piers only instead of immensely thick walls as in the tomb; but any thrust which might exist was counteracted by its transmission across aisles to the outer wall.
Islamic architecture:
The Muqarnas is a form of vaulting common in Islamic architecture.
Modern vaults:
Hyperbolic paraboloids The 20th century saw great advances in reinforced concrete design. The advent of shell construction and the better mathematical understanding of hyperbolic paraboloids allowed very thin, strong vaults to be constructed with previously unseen shapes. The vaults in the Church of Saint Sava are made of prefabricated concrete boxes. They were built on the ground and lifted to 40 m on chains.
Modern vaults:
Vegetal vault When made by plants or trees, either artificially or grown on purpose by humans, structures of this type are called tree tunnels. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Renfe Class 446**
Renfe Class 446:
The Renfe Class 446 is a series of electric multiple units designed to provide effective commuter services in major urban centers in Spain.
History:
Planning began in the early 1980s and resulted in the prototype Class 445 when commuter services began to assume a major role in cities such as Madrid. The trains then operating these services, mainly Class 440 units, were becoming inadequate to cope with the volume of traffic. It was decided that to provide the required facilities for these service, with their short inter-station distances and high passenger volume, new trains would be needed, with greater acceleration at the cost of maximum speed. Due to the 2004 Madrid train bombings, three trains of this class were decommissioned. Two of the damaged sets were amalgamated to form a new train and the third set was repaired and returned to service.
Design:
Introduced in 1989, the 446 series was an innovative train design by Renfe. Each set having 2,400 kW (3,200 hp) of power, allowing a maximum acceleration of 1 m (3 ft 3 in) /s² with a top speed of 100 km/h (62 mph). The three car sets consist of two driving motor cars and an intermediate trailer, a new concept for Renfe. Each car has three pairs of double doors for quick entry and exit of travelers. One of the major drawbacks of these units, which earned them the nickname dodotis (a type of baby's nappy), is that originally they did not have toilets in the cars.
Design:
Outwardly they are of similar appearance to the later Renfe Class 447 design, with which they can operate with some limitations in traction and brake. Recent changes and enhancements included the addition of hustle alarms, both sound and indicator lights, for closing doors. Of the 170 units delivered, 25 have been withdrawn, including those affected by the 2004 Madrid train bombings. The units can be found operating metro type inner suburban services around Madrid, Seville, San Sebastián, Santander and Bilbao where the distances between stations is often 1–2 km (0.62–1.24 mi) or even less.
Formation:
Each set is made of 3 cars: Numbering example ^ The set no.1 is numbered as 001M-001R-002M (UIC numbering: 9-446-001+7-446-001+9-446-002), the second as 003M-002R-004M (9-446-003+7-446-002+9-446-004), and so on (M means powered and R trailer). | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**PCDHA5**
PCDHA5:
Protocadherin alpha-5 is a protein that in humans is encoded by the PCDHA5 gene.This gene is a member of the protocadherin alpha gene cluster, one of three related gene clusters tandemly linked on chromosome 5 that demonstrate an unusual genomic organization similar to that of B-cell and T-cell receptor gene clusters. The alpha gene cluster is composed of 15 cadherin superfamily genes related to the mouse CNR genes and consists of 13 highly similar and 2 more distantly related coding sequences. The tandem array of 15 N-terminal exons, or variable exons, are followed by downstream C-terminal exons, or constant exons, which are shared by all genes in the cluster. The large, uninterrupted N-terminal exons each encode six cadherin ectodomains while the C-terminal exons encode the cytoplasmic domain. These neural cadherin-like cell adhesion proteins are integral plasma membrane proteins that most likely play a critical role in the establishment and function of specific cell-cell connections in the brain. Alternative splicing has been observed and additional variants have been suggested but their full-length nature has yet to be determined. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Sequence space**
Sequence space:
In functional analysis and related areas of mathematics, a sequence space is a vector space whose elements are infinite sequences of real or complex numbers. Equivalently, it is a function space whose elements are functions from the natural numbers to the field K of real or complex numbers. The set of all such functions is naturally identified with the set of all possible infinite sequences with elements in K, and can be turned into a vector space under the operations of pointwise addition of functions and pointwise scalar multiplication. All sequence spaces are linear subspaces of this space. Sequence spaces are typically equipped with a norm, or at least the structure of a topological vector space.
Sequence space:
The most important sequence spaces in analysis are the ℓp spaces, consisting of the p-power summable sequences, with the p-norm. These are special cases of Lp spaces for the counting measure on the set of natural numbers. Other important classes of sequences like convergent sequences or null sequences form sequence spaces, respectively denoted c and c0, with the sup norm. Any sequence space can also be equipped with the topology of pointwise convergence, under which it becomes a special kind of Fréchet space called FK-space.
Definition:
A sequence x∙=(xn)n∈N in a set X is just an X -valued map x∙:N→X whose value at n∈N is denoted by xn instead of the usual parentheses notation x(n).
Space of all sequences Let K denote the field either of real or complex numbers. The set KN of all sequences of elements of K is a vector space for componentwise addition (xn)n∈N+(yn)n∈N=(xn+yn)n∈N, and componentwise scalar multiplication α(xn)n∈N=(αxn)n∈N.
A sequence space is any linear subspace of KN.
Definition:
As a topological space, KN is naturally endowed with the product topology. Under this topology, KN is Fréchet, meaning that it is a complete, metrizable, locally convex topological vector space (TVS). However, this topology is rather pathological: there are no continuous norms on KN (and thus the product topology cannot be defined by any norm). Among Fréchet spaces, KN is minimal in having no continuous norms: But the product topology is also unavoidable: KN does not admit a strictly coarser Hausdorff, locally convex topology. For that reason, the study of sequences begins by finding a strict linear subspace of interest, and endowing it with a topology different from the subspace topology.
Definition:
ℓp spaces For 0<p<∞, ℓp is the subspace of KN consisting of all sequences x∙=(xn)n∈N satisfying If p≥1, then the real-valued function ‖⋅‖p on ℓp defined by defines a norm on ℓp.
In fact, ℓp is a complete metric space with respect to this norm, and therefore is a Banach space.
If p=2 then ℓ2 is also a Hilbert space when endowed with its canonical inner product, called the Euclidean inner product, defined for all x∙,y∙∈ℓp by The canonical norm induced by this inner product is the usual ℓ2 -norm, meaning that ‖x‖2=⟨x,x⟩ for all x∈ℓp.
If p=∞, then ℓ∞ is defined to be the space of all bounded sequences endowed with the norm ℓ∞ is also a Banach space.
Definition:
If 0<p<1, then ℓp does not carry a norm, but rather a metric defined by c, c0 and c00 A convergent sequence is any sequence x∙∈KN such that lim n→∞xn exists. The set c of all convergent sequences is a vector subspace of KN called the space of convergent sequences. Since every convergent sequence is bounded, c is a linear subspace of ℓ∞.
Definition:
Moreover, this sequence space is a closed subspace of ℓ∞ with respect to the supremum norm, and so it is a Banach space with respect to this norm.
Definition:
A sequence that converges to 0 is called a null sequence and is said to vanish. The set of all sequences that converge to 0 is a closed vector subspace of c that when endowed with the supremum norm becomes a Banach space that is denoted by c0 and is called the space of null sequences or the space of vanishing sequences. The space of eventually zero sequences, 00 , is the subspace of c0 consisting of all sequences which have only finitely many nonzero elements. This is not a closed subspace and therefore is not a Banach space with respect to the infinity norm. For example, the sequence (xnk)k∈N where xnk=1/k for the first n entries (for k=1,…,n ) and is zero everywhere else (that is, (xnk)k∈N=(1,1/2,…,1/(n−1),1/n,0,0,…) ) is a Cauchy sequence but it does not converge to a sequence in 00 .
Definition:
Space of all finite sequences Let all but finitely many equal 0} ,denote the space of finite sequences over K . As a vector space, K∞ is equal to 00 , but K∞ has a different topology. For every natural number n∈N , let Kn denote the usual Euclidean space endowed with the Euclidean topology and let In Kn:Kn→K∞ denote the canonical inclusion In Kn(x1,…,xn)=(x1,…,xn,0,0,…) .The image of each inclusion is Im In Kn)={(x1,…,xn,0,0,…):x1,…,xn∈K}=Kn×{(0,0,…)} and consequently, Im In Kn).
Definition:
This family of inclusions gives K∞ a final topology τ∞ , defined to be the finest topology on K∞ such that all the inclusions are continuous (an example of a coherent topology). With this topology, K∞ becomes a complete, Hausdorff, locally convex, sequential, topological vector space that is not Fréchet–Urysohn. The topology τ∞ is also strictly finer than the subspace topology induced on K∞ by KN . Convergence in τ∞ has a natural description: if v∈K∞ and v∙ is a sequence in K∞ then v∙→v in τ∞ if and only v∙ is eventually contained in a single image Im In Kn) and v∙→v under the natural topology of that image. Often, each image Im In Kn) is identified with the corresponding Kn ; explicitly, the elements (x1,…,xn)∈Kn and (x1,…,xn,0,0,0,…) are identified. This is facilitated by the fact that the subspace topology on Im In Kn) , the quotient topology from the map In Kn , and the Euclidean topology on Kn all coincide. With this identification, In Kn)n∈N) is the direct limit of the directed system In Km→Kn)m≤n∈N,N), where every inclusion adds trailing zeros: In Km→Kn(x1,…,xm)=(x1,…,xm,0,…,0) .This shows (K∞,τ∞) is an LB-space.
Definition:
Other sequence spaces The space of bounded series, denote by bs, is the space of sequences x for which sup n|∑i=0nxi|<∞.
This space, when equipped with the norm sup n|∑i=0nxi|, is a Banach space isometrically isomorphic to ℓ∞, via the linear mapping (xn)n∈N↦(∑i=0nxi)n∈N.
The subspace cs consisting of all convergent series is a subspace that goes over to the space c under this isomorphism.
The space Φ or 00 is defined to be the space of all infinite sequences with only a finite number of non-zero terms (sequences with finite support). This set is dense in many sequence spaces.
Properties of ℓp spaces and the space c0:
The space ℓ2 is the only ℓp space that is a Hilbert space, since any norm that is induced by an inner product should satisfy the parallelogram law ‖x+y‖p2+‖x−y‖p2=2‖x‖p2+2‖y‖p2.
Substituting two distinct unit vectors for x and y directly shows that the identity is not true unless p = 2.
Properties of ℓp spaces and the space c0:
Each ℓp is distinct, in that ℓp is a strict subset of ℓs whenever p < s; furthermore, ℓp is not linearly isomorphic to ℓs when p ≠ s. In fact, by Pitt's theorem (Pitt 1936), every bounded linear operator from ℓs to ℓp is compact when p < s. No such operator can be an isomorphism; and further, it cannot be an isomorphism on any infinite-dimensional subspace of ℓs, and is thus said to be strictly singular.
Properties of ℓp spaces and the space c0:
If 1 < p < ∞, then the (continuous) dual space of ℓp is isometrically isomorphic to ℓq, where q is the Hölder conjugate of p: 1/p + 1/q = 1. The specific isomorphism associates to an element x of ℓq the functional for y in ℓp. Hölder's inequality implies that Lx is a bounded linear functional on ℓp, and in fact so that the operator norm satisfies sup y∈ℓp,y≠0|Lx(y)|‖y‖p≤‖x‖q.
Properties of ℓp spaces and the space c0:
In fact, taking y to be the element of ℓp with if if xn≠0 gives Lx(y) = ||x||q, so that in fact ‖Lx‖(ℓp)∗=‖x‖q.
Properties of ℓp spaces and the space c0:
Conversely, given a bounded linear functional L on ℓp, the sequence defined by xn = L(en) lies in ℓq. Thus the mapping x↦Lx gives an isometry The map ℓq→κq(ℓp)∗→(κq∗)−1 obtained by composing κp with the inverse of its transpose coincides with the canonical injection of ℓq into its double dual. As a consequence ℓq is a reflexive space. By abuse of notation, it is typical to identify ℓq with the dual of ℓp: (ℓp)* = ℓq. Then reflexivity is understood by the sequence of identifications (ℓp)** = (ℓq)* = ℓp.
Properties of ℓp spaces and the space c0:
The space c0 is defined as the space of all sequences converging to zero, with norm identical to ||x||∞. It is a closed subspace of ℓ∞, hence a Banach space. The dual of c0 is ℓ1; the dual of ℓ1 is ℓ∞. For the case of natural numbers index set, the ℓp and c0 are separable, with the sole exception of ℓ∞. The dual of ℓ∞ is the ba space.
Properties of ℓp spaces and the space c0:
The spaces c0 and ℓp (for 1 ≤ p < ∞) have a canonical unconditional Schauder basis {ei | i = 1, 2,...}, where ei is the sequence which is zero but for a 1 in the i th entry.
The space ℓ1 has the Schur property: In ℓ1, any sequence that is weakly convergent is also strongly convergent (Schur 1921). However, since the weak topology on infinite-dimensional spaces is strictly weaker than the strong topology, there are nets in ℓ1 that are weak convergent but not strong convergent.
Properties of ℓp spaces and the space c0:
The ℓp spaces can be embedded into many Banach spaces. The question of whether every infinite-dimensional Banach space contains an isomorph of some ℓp or of c0, was answered negatively by B. S. Tsirelson's construction of Tsirelson space in 1974. The dual statement, that every separable Banach space is linearly isometric to a quotient space of ℓ1, was answered in the affirmative by Banach & Mazur (1933). That is, for every separable Banach space X, there exists a quotient map Q:ℓ1→X , so that X is isomorphic to ker Q . In general, ker Q is not complemented in ℓ1, that is, there does not exist a subspace Y of ℓ1 such that ker Q . In fact, ℓ1 has uncountably many uncomplemented subspaces that are not isomorphic to one another (for example, take X=ℓp ; since there are uncountably many such X's, and since no ℓp is isomorphic to any other, there are thus uncountably many ker Q's).
Properties of ℓp spaces and the space c0:
Except for the trivial finite-dimensional case, an unusual feature of ℓp is that it is not polynomially reflexive.
Properties of ℓp spaces and the space c0:
ℓp spaces are increasing in p For p∈[1,∞] , the spaces ℓp are increasing in p , with the inclusion operator being continuous: for 1≤p<q≤∞ , one has ‖x‖q≤‖x‖p . Indeed, the inequality is homogeneous in the xi , so it is sufficient to prove it under the assumption that ‖x‖p=1 . In this case, we need only show that ∑|xi|q≤1 for q>p . But if ‖x‖p=1 , then |xi|≤1 for all i , and then ∑|xi|q≤∑|xi|p=1 ℓ2 is isomorphic to all separable, infinite dimensional Hilbert spaces Let H be a separable Hilbert space. Every orthogonal set in H is at most countable (i.e. has finite dimension or ℵ0 ). The following two items are related: If H is infinite dimensional, then it is isomorphic to ℓ2 If dim(H) = N, then H is isomorphic to CN
Properties of ℓ1 spaces:
A sequence of elements in ℓ1 converges in the space of complex sequences ℓ1 if and only if it converges weakly in this space. If K is a subset of this space, then the following are equivalent: K is compact; K is weakly compact; K is bounded, closed, and equismall at infinity.Here K being equismall at infinity means that for every ε>0 , there exists a natural number nε≥0 such that {\textstyle \sum _{n=n_{\epsilon }}^{\infty }|s_{n}|<\varepsilon } for all s=(sn)n=1∞∈K | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**Baer function**
Baer function:
Baer functions Bpq(z) and Cpq(z) , named after Karl Baer, are solutions of the Baer differential equation d2Bdz2+12[1z−b+1z−c]dBdz−[p(p+1)z+q(b+c)(z−b)(z−c)]B=0 which arises when separation of variables is applied to the Laplace equation in paraboloidal coordinates. The Baer functions are defined as the series solutions about z=0 which satisfy Bpq(0)=0 , Cpq(0)=1 . By substituting a power series Ansatz into the differential equation, formal series can be constructed for the Baer functions. For special values of p and q , simpler solutions may exist. For instance, ln [z+(z−b)(z−c)−(b+c)/2bc−(b+c)/2] Moreover, Mathieu functions are special-case solutions of the Baer equation, since the latter reduces to the Mathieu differential equation when b=0 and c=1 , and making the change of variable cos 2t Like the Mathieu differential equation, the Baer equation has two regular singular points (at z=b and z=c ), and one irregular singular point at infinity. Thus, in contrast with many other special functions of mathematical physics, Baer functions cannot in general be expressed in terms of hypergeometric functions.
Baer function:
The Baer wave equation is a generalization which results from separating variables in the Helmholtz equation in paraboloidal coordinates: d2Bdz2+12[1z−b+1z−c]dBdz+[k2z2−p(p+1)z−q(b+c)(z−b)(z−c)]B=0 which reduces to the original Baer equation when k=0 | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
**27-Hydroxycholesterol**
27-Hydroxycholesterol:
27-Hydroxycholesterol (27-HC) is an endogenous oxysterol with multiple biological functions, including activity as a selective estrogen receptor modulator (SERM) (a mixed, tissue-specific agonist-antagonist of the estrogen receptor (ER)) and as an agonist of the liver X receptor (LXR). It is a metabolite of cholesterol that is produced by the enzyme CYP27A1.A link between high cholesterol and breast cancer has been identified, and it has been proposed that this is due to 27-HC production by CYP27A1. Because of its estrogenic action, 27-HC stimulates the growth of ER-positive breast cancer cells, and has been implicated in limiting the effectiveness of aromatase inhibitors in the treatment of breast cancer. As such, identified CYP27A1 inhibitors, including the marketed drugs anastrozole, fadrozole, bicalutamide, dexmedetomidine, ravuconazole, and posaconazole, have been proposed as potential adjuvant therapies in ER-positive breast cancer. | kaggle.com/datasets/mbanaei/all-paraphs-parsed-expanded |
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