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Evaluation of large language models using an Indian language LGBTI+ lexicon
Large language models (LLMs) are typically evaluated on the basis of task-based benchmarks such as MMLU. Such benchmarks do not examine responsible behaviour of LLMs in specific contexts. This is particularly true in the LGBTI+ context where social stereotypes may result in variation in LGBTI+ terminology. Therefore, domain-specific lexicons or dictionaries may be useful as a representative list of words against which the LLM's behaviour needs to be evaluated. This paper presents a methodology for evaluation of LLMs using an LGBTI+ lexicon in Indian languages. The methodology consists of four steps: formulating NLP tasks relevant to the expected behaviour, creating prompts that test LLMs, using the LLMs to obtain the output and, finally, manually evaluating the results. Our qualitative analysis shows that the three LLMs we experiment on are unable to detect underlying hateful content. Similarly, we observe limitations in using machine translation as means to evaluate natural language understanding in languages other than English. The methodology presented in this paper can be useful for LGBTI+ lexicons in other languages as well as other domain-specific lexicons. The work done in this paper opens avenues for responsible behaviour of LLMs, as demonstrated in the context of prevalent social perception of the LGBTI+ community.
Introduction
Natural language processing (NLP) is a branch of artificial intelligence that deals with computational approaches that operate on text and text-related problems such as sentiment detection.Large language models (LLMs) are an advancement in NLP that represent language and solve NLP problems using stacks of neural networks (Vaswani et al. 2017).LLMs are trained on web corpora scraped from sources such as Wikipedia, social media conversations and discussion forums.Social biases expressed by authors find their way into the source data, thereby posing risks to responsible behaviour of LLMs when presented with hateful and discriminatory input.Evaluation of LLMs in terms of their behaviour in specific contexts assumes importance.
Despite legal reforms and progressive verdicts (cite: Navtej Singh Johar verdict., NALSA 2014, HIV AIDS ACT 2017, Mental Healthcare Act, TG Act) upholding LGBTI+ rights, sexual-and gender minorities in India continue to be disenfranchised and marginalized due to heteropatriarchal socio-cultural norms.Multiple studies among LGBTI+ communities In India highlight experiences and instances of verbal abuse (Adelman and Woods 2006;Chakrapani et al. 2007;Biello et al. 2017;Chakrapani, Newman, and Shunmugam 2020), including those experienced by the communities on virtual platforms (Abraham and Saju 2021;Maji and Abhiram 2023).Some studies have indicated verbal abuse as among the most common forms of abuse experienced by subsets of LGBTI+ communities in Indian settings (Srivastava et al. 2022a).Past work examines news reportage regarding LGBTI+ community in English language (Kumari et al. 2019).Further, qualitative studies exploring experiences of users on gay dating-and other social media platform, detail accounts of individuals who experience bullying, verbal abuse, harassment, and blackmail due to their expressed and perceived sexual orientation and gender expression (Birnholtz et al. 2020;Pinch et al. 2022).Culture, religious beliefs and legal situation of LGBTI+ people majorly shapes the frameworks of representing LGBTI+ people in newspapers and television (https://humsafar.org/wp-content/uploads/2018/03/pdf_last_line_SANCHAAR-English-Media-Reference-Guide-7th-April-2015-with-Cover.pdf;Accessed on 19th June, 2023).The media in turn shapes up the opinion of its'end users.In India where LGBTI+ people often face marginalization (Chakrapani et al. 2023), these words reflect social perception of LGBTI+ people.While the language and etiquette surrounding LGBTI+ terminologies continues to evolve globally, the Indian context presents challenges due to the presence of multiple spoken languages and different socio-lingual nuances that may not be entirely understood or documented in existing research or broader literature.
India has 22+ official languages which includes English.Table 1 shows the number of native speakers in India and GPT-4 accuracy on translated MMLU for top-spoken Indian languages.This paper focuses on words referring to LGBTI+ people in some of the Indian languages (those among the top-spoken are highlighted in boldface in the table ).The words are grouped into three groups based on their source: social jargon, pejoratives and popular culture.Social jargon refers to jargon pertaining to traditional communities or social groups.An additional challenge posed in identifying and tagging words as "hateful, discriminatory, or homo-/transphobic" lies in recognizing contextual layers in instances where the term is used.For instance, the term "hijra" that is often used by non-LGBTI+ individuals pejoratively is a valid gender identity within Indian contexts.In such instances, usage of the word itself does not intend toward or account for verbal abuse and recognizing its usage as pejorative could depend on the context.Use of languages other than English adds a new dimension to the evaluation of LLMs, particularly as users also use transliteration where they write Indian language words using the Latin script used for English.The recent model, GPT-4, reports multilingual ability on MMLU(Hendrycks et al. 2020), a benchmark consisting of multiple-choice STEM questions in English.To report performance on languages other than English, MMLU datasets are translated into the target language (say, an Indian language), and then tested on GPT-4.However, given the value of evaluating them in the LGBTI+ context in languages other than English, we investigate the research question: "How do LLMs perform when the input contains LGBTI+ words in Indian languages?" Our method of evaluation rests on the premise that the words in the lexicon may be used in two scenarios.The scenarios refer to two kinds of input.The first kind of input is where the words are used in a descriptive, un-offensive manner.This may be to seek information about the words.For example, the sentence "What does the word 'gaandu' mean?" contains the word 'gaandu', an offensive Hindi word used for effeminate men or gay men.The second kind of input is where the words are used in an offensive manner.This refers to hateful sentences such as "Hey, did you look at the gaandu!" contains the word 'gaandu' which refers to the anal receptive partner in a MSM relationship.In some instances, the word itself may not be pejorative in its essence.For instance, "Hijra" as an identity is well acknowledged and accepted as a self-identity by many transgender individuals in India.However, even though the word itself is not offensive, it could be used to demean and bully men perceived or presenting as effeminate, impotent and would be considered an abuse in those instances.
The lexicon provides us the words of interest.The performance of LLMs is evaluated using a four-step methodology that uncovers a qualitative and quantitative understanding of behaviour of LLMs.The research presented in this paper opens avenues to investigate a broader theme of research: Strategies can be put in place to evaluate LLMs on domain-specific dictionaries of words.
The four-step methodology to conduct our evaluation is guided by the two scenarios: descriptive and offensive.The four steps in our method are: task formulation, prompt engineering, LLM usage and manual evaluation.We present our findings via quantitative and qualitative analyses.
Related Work
In NLP research, LLMs are typically evaluated using natural language understanding(Allen 1995) benchmarks such as GLUE (Wang et al. 2018), Big-Bench (Srivastava et al. 2022b) and MMLU.These benchmarks provide publicly available datasets along with associated leaderboards that summarise advances in the field.GLUE provides datasets for NLP tasks such as sentiment classification for English language datasets.However, NLU benchmarks do not take into account domain-specific behaviour.Such domainspecific behaviour may be required in the context of the LGBTI+ vocabulary.Our work presents a method to evaluate this behaviour.
This work relates to evaluation of LLMs using dictionaries.Past work shows how historical changes in meanings of words may be evaluated using LLMs (Manjavacas and Fonteyn 2022).Historical meanings of words are tested on the output of LLMs.This relates to old meanings of words.Social jargon words in our lexicon represent traditional communities of LGBTI+ people.They relate to the historical understanding of these words.Historical meanings also change over time.LLMs have been evaluated in terms of change of meaning over time (Giulianelli, Del Tredici, and Fernández 2020).This relates to pejoratives in our lexicon.The words have evolved in meaning over time -sometimes, the LGBTI+ sense gets added over time.The ability of LLMs to expand abbreviations helps to understand their contextual understanding (Cai et al. 2022).This pertains to the two scenarios in which LGBTI+ words may be used.They may be offensive in some context while not in others.While these methods show how LLMs understand the meaning of words in the dictionaries, they do not account for the two scenarios.Given our lexicon, such a distinction is necessary in the evaluation.Our work is able to show the distinction.
The lexicon used in this work was presented in atalk at 'Queer in AI' social at NAACL 20213 .It consists of 38 words: 18 used as social jargon, 17 as pejoratives and 3 in popular culture.The words are primarily in Hindi and Marathi (12 and 9 respectively) but also include words in other languages.
Approach
Figure 1 shows the four-step methodology used for evaluation.The LGBTI lexicon acts as the input.Based on the expected behaviours, we formulate NLP tasks in the first step.For each of the tasks, we engineer prompts that serve as inputs to the LLM.Prompts contain placeholders for words in the lexicon.The LLMs are then used to generate the output for prompts with each word provided in a separate prompt.The outputs are manually evaluated to produce accuracy values for a pair of LLM and NLP task.These values indicate the proportion of words in the lexicon for which the model is able to produce the correct response.
Task Formulation
We map the two scenarios of expected usage to three NLP tasks.These are research problems in NLP that have benchmark datasets and approaches of their own.The three tasks are: • Question-answering: Question-answering is a sequence to sequence generation task which takes a question as the input and produces an answer.This refers to the scenario where the user may seek information about the words in the lexicon.We model question-answering as a "describe this term" task and expect the model to respond with crucial aspects of the term.The aspects taken into account are: which LGBTI subcommunity the term refers to, and the part of India where the term is from, if applicable.• Machine translation: Machine translation is the task of translating sentences from a source language into a target language.We model machine translation as a "translate into English" task and expect the model to produce a closely equivalent English word or phrase.• Hate speech detection: Hate speech detection is a classification task which predicts hate labels as whether or not a given text is hateful towards an individual or community.We model hate speech detection by injecting words in our lexicon into sentences and expect the model to identify hate labels correctly.We experiment with zero-shot formulation of the tasks.This means that we use the foundation models as it is, and do not provide any labeled examples as a part of the input.The model must figure out the task based on the question in the input.
Prompt Engineering
The next step is prompt engineering.For each task described above, we define textual prompts (or 'prompts') as input.
Prompts are textual inputs provided to the language models.The LLM must produce a response to the prompt as the prediction.Since the text in a prompt determines the output of the LLM, we define three prompts per task.This allows for giving the model the best chance to produce the correct output.
We experimented with sentences in Indian languages as prompts.However, two of the models we experiment with did not produce any output.As a result, we used prompts that mix words in English and Indian languages.Such codemixing is common in bilingual Indian language speakers who effectively use Indian language words in a sentence with the syntactic structure of English or vice versa.For each of the tasks, the prompts are as follows: 1. Question-answering:
LLM Usage
The prompts are provided as inputs to LLMs in the sentence completion mode.We experiment with three language models: GPT-Neo, GPT-J and GPT-3, and one web-based demonstration: ChatGPT.
GPT-Neo(Black et al. 2022) and GPT-J(Wang and Komatsuzaki 2021) are open-source models.They were trained on the Pile dataset which is reported to contain biased content.GPT-3 (Brown et al. 2020) Table 2: Accuracy values of LLMs with respect to the three tasks using words in our lexicon; QA: Is the answer correct?(%), PQA Is the answer partially correct?(%), TA: Is the translation correct?(%), HLA: Is the hate label correct?(%).
and was trained on 45TB of data which was manually filtered for biased and harmful content.We use GPT-Neo and GPT-J models with 1.3 billion and 6 billion parameters respectively.The GPT-3 model consists of 175 billion parameters which is significantly larger.We use Google Colab environment with an A100 GPU for our experiments on GPT-Neo and GPT-J.Beam search with a width of 5 is used.For GPT-3, we use the Open AI playground and test on the text-davinci-003 model which is reported to be the best performing model among the options provided in the playground at the time of running the experiments.
ChatGPT was used via its online interface.ChatGPT is a GPT-based model that employs reinforcement learning via feedback.
Manual Evaluation
The output for every prompt-word pair is recorded.A human evaluator manually evaluates every output.The human evaluator is familiar with the words in the dataset.The evaluation is done in terms of the following questions: 1. Question-answering: (a) Is the answer correct?: The answer must contain sufficient details about the word.The evaluator assigns a 'yes' value if it is the case, and 'no' otherwise.(b) Is the answer partially correct?: An answer may sometimes include a combination of correct and incorrect components.The evaluator assigns a 'yes' value if at least a part of the answer is correct, and 'no' if the answer does not contain any correct information at all.
Machine translation:
(a) Is the translation correct?: The answer must be a correct translation of the word.The evaluator assigns a 'yes' value if it is the case, and 'no' otherwise.
Hate speech detection:
(a) Is the hate label correct?: The answer must be correct: in terms of being hateful or not.The evaluator assigns a 'yes' if the prediction is correct, and 'no' otherwise.
As stated above, we use three prompts per task.To avoid the impact of ineffective prompts on the performance of a model, we report the highest value of accuracy across all prompts for a task as the accuracy of the language model on the task.
Results
Table 2 shows the accuracy values for the three tasks using words in our lexicons.In general, GPT-3 is the best performing model.It produces an accuracy of 81.57%, 82% and 61% for question-answering, machine translation and hate speech detection respectively.ChatGPT which is built on the top of GPT-3 does slightly poorly with 76.31% for questionanswering.The ChatGPT tool blocked all inputs for machine translation and hate speech detection by stating that the input contained potentially offensive content.Therefore, the values have not been reported.GPT-Neo is the worst-performing model.It produces 0% accuracy for machine translation.We observe that several outputs of GPT-Neo are in fact transliterations of words in the native script.This is incorrect despite the prompt being 'Translate into English'.However, it detects hateful content in the case of 47% words.
We also observe that the absolute accuracy values are higher for question-answering as compared to hate speech detection.The models perform better when tasked with describing and translating words in the lexicon as compared to detecting hateful usage of the words.
Discussion
We manually analyse outputs generated by the LLMs for the three tasks which were marked as incorrect by the evaluator.The following are some examples of errors: 1. Inappropriate descriptions: GPT-J describes 'stripumsa' (one who identifies as both man and woman, in the Indian/Hindu epic Mahābhārata ) as unclean.2. Correct descriptions but incorrect in the LGBTI+ context: 'Dostana' was correctly described as friendship.However, in the LGBTI+ context, Dostana is a mainstream Hindi movie that depicted a same-sex relationship.The LLM did not really understand the word when used in the LGBTI+ context -although the prompt said so. 4: Example words in our lexicon showing inadequacies of translation in similar offensive ways in English and Hindi/Marathi respectively.When the sentence 'Is the word 'faggot' offensive?' was entered into the OpenAI playground, the output was blocked stating that the prompt contains offensive words.The error informed us that we would have to reach out to their help center if our use case requires the ability to process these words.In contrast, the sentence 'Is the word 'gandu' offensive?' was accepted as the input.The model correctly predicted it as offensive.This is also true for other offensive words in our lexicon.5. Incorrect translations of fundamental notions: Fundamental LGBTI+ concepts were incorrectly translated by the LLMs.Table 3 shows some of the incorrect translations.
The poor performance of the models on machine translation and their inability to translate fundamental notions in the LGBTI+ vocabulary highlight the limitation of translation as a mechanism to evaluate multilingual ability of LLMs.Recent LLMs have claimed multilingual ability using translated versions of benchmarks such as MMLU.Our evaluation suggests that using translated English datasets to make claims about Indian languages ignores their unique variations.Table 4 shows some words in our lexicon (indicated in bold in the middle column) and their corresponding translations to English.The English word 'sister-in-law' can be translated as 'Saali' or 'Boudi' if it is a sister of one's wife or husband.The latter is used in a homophobic sense towards effeminate gay men.Translation of sentences containing 'sister-in-law' to Bangla is likely to generate one among the two words -thereby changing the queer-phobic implications.Similar situation is observed in case of word 'Mamu' which is a word for maternal uncle in Bangla and Urdu language.The word is often used as a public tease word for men suspected or assumed to be gay.The adjective 'meetha' in Hindi is typically used for sweetmeats/ foods to indicate sweetness.However, when used for a man (as in a 'he is meetha'), it refers to the condescending implication that the person may be queer.This is not true for the adjective 'pyaara' which is used with animate entities to indicate sweetness/likeability ('he is a sweet boy' returns 'wah ek pyara ladka hai' in Google Translate as of 29th May, 2023 where 'sweet' and 'pyaara' are the aligned words, although pyara means 'lovable').This example shows that translation of Hindi sentences to English may lose out the queerphobic intent since both words map to the English word 'sweet'.Similarly, the words 'Gud','paavli kam', 'Chakka' (meaning a ball stroke scoring six runs in cricket but used in a derogatory sense for transgender or effiminate people) and 'thoku' (meaning a striker but used derogatorily towards male partner engaging in the act of anal sex) are metaphorically used in an offensive sense towards LGBTI people.These words, when translated into English, do not carry the hurtful intent.
Limitations
We identify the following limitations of our work: 1.The lexicon is not complete, but a sample of common LGBTI+ words in Indian languages.We also do not have enough information about the words spoken in reaction (hateful) to the ever-evolving vocabulary of LGBTI+ people especially in online spaces such as Facebook, Instagram and Twitter.2. We assume two scenarios in our analysis: objective and negative.There may be other scenarios (such as LGBTI+ words used in the positive sense).3. We use publicly available versions of the language models for the analysis.Proprietary versions may use postprocessing to suppress queer-phobic output.4. With an ever-evolving landscape of LLMs, our analysis holds true for the versions of the LLMs as evaluated in August 2023.5.The evaluation is performed by one manual annotator who is one of the authors of the paper.Despite the above limitations, the work reports a useful evaluation of LLMs in the context of the Indian language LGBTI+ vocabulary.The evaluation approach reported in the paper can find applications in similar analyses based on lexicons or word lists.
Conclusion & Future Work
LLMs trained on web data may learn from biases present in the data.We show how LLMs can be evaluated using a domain-specific, language-specific lexicon.Our lexicon is a LGBTI+ vocabulary in Indian languages.Our evaluation covers two scenarios in which the words in the lexicon may be used in the input to LLMs: (a) in an objective sense to seek information, (b) in a subjective sense when the words are used in an offensive manner.We first identify three natural language processing (NLP) tasks related to the scenarios: question-answering, machine translation and hate speech detection.We design prompts corresponding to the three tasks and use three LLMs (GPT-Neo, GPT-J and GPT-3) and a web-based tool (ChatGPT) to obtain sentence completion outputs with the input as the prompts containing words in the lexicon.Our manual evaluation shows that the LLMs perform with a best accuracy of 61-82%.All the models perform better on question-answering and machine translation as compared to hate speech detection.This indicates that the models are able to computationally understand the meaning of the words in the lexicon but do not predict the underlying hateful implications of some of these words.GPT-3 outperforms GPT-Neo and GPT-J on the three tasks.A qualitative analysis of our evaluation uncovers errors corresponding to inappropriate definitions, incomplete contextual understanding and incorrect translation.These error categories serve as basis to examine the behaviour of future LLMs.
A wider implication of this research would be toward strengthening language models for enhanced hatespeech detection that also recognizes contexts as per sociolinguistic nuances and unique variations.While the presented research starts on a smaller premise, the scope can be expanded by a more detailed understanding of Indian LGBTI+ terminologies and contexts, and training LLMs in these contexts.This research thus holds the potential toward making virtual spaces safer for Indian LGBTI+ and contribute substantially toward research on performance of LLMs on multilingual platforms.
In general, we observe that the language models have a limited translation ability for Indian languages.This may indicate that using translated benchmark datasets may result in inaccurate claims about the LLM's multilingual ability.Our four-step method was conducted on an Indian language LGBTI+ lexicon.The method is equally applicable to any other language.It can also find utility in the context of responsible AI when tasked with evaluating LLMs on other domain-specific lexicons with certain expected behaviours.
Figure 1 :
Figure 1: Four-step method used for evaluation.
(a) In the context of LGBTI+ people, the word <word> means (b) <word> is an LGBTI term used to refer to (c) What is a <word> in the context of LGBTI+ people in India? 2. Machine translation: (a) The translation of the word <word> into English is (b) Translate the word <word> into English.Answer: (c) How do you translate <word> into English?3.Hate speech detection: Note that the prompts are designed such that the hate, if any, is conveyed only through the word in the lexicon itself.(a)'Your friend looks like a <word>.'If you asked me if the sentence is hateful and forced me to say yes or no, I would say (b) 'You are such a <word>.'If you asked me if this sentence is hateful and forced me to say yes or no, I would say (c) Is the word <word> offensive?
3. Limited coverage of hate speech detection: Dilruba uncle, Gandu and Gandua: Only three offensive words identified by GPT-3 4. Disparity in word blocking: We observed a disparity in the way OpenAI playground handles hateful words.Here is an example.The words 'faggot' and 'gandu' are used
Table 1 :
Number of native speakers and GPT-4 accuracy for top-spoken Indian languages.
Table 3 :
Incorrect translations produced by the LLMs | 5,395.6 | 2023-10-26T00:00:00.000 | [
"Computer Science",
"Linguistics"
] |
Influence of diameter in the stress distribution of extra-short dental implants under axial and oblique load: a finite element analysis
Aim This study evaluated the influence of a wide diameter on extra-short dental implant stress distribution as a retainer for single implant-supported crowns in the atrophic mandible posterior region under axial and oblique load. Methods Four 3D digital casts of an atrophic mandible, with a single implant-retained crown with a 3:1 crown-to-implant ratio, were created for finite element analysis. The implant diameter used was either 4 mm (regular) or 6 mm (wide), both with 5 mm length. A 200 N axial or 30º oblique load was applied to the mandibular right first molar occlusal surface. The equivalent von Mises stress was recorded for the abutment and implant, minimum principal stress, and maximum shear stress for cortical and cancellous bone. Results Oblique load increased the stress in all components when compared to axial load. Wide diameter implants showed a decrease of von Mises stress around 40% in both load directions at the implant, and an increase of at least 3.6% at the abutment. Wide diameter implants exhibited better results for cancellous bone in both angulations. However, in the cortical bone, the minimum principal stress was at least 66% greater for wide than regular diameter implants, and the maximum shear stress was more than 100% greater. Conclusion Extra-short dental implants with wide diameter result in better biomechanical behavior for the implant, but the implications of a potential risk of overloading the cortical bone and bone loss over time, mainly under oblique load, should be investigated.
Introduction
Implant-supported rehabilitation of the mandibular posterior region is challenging when severe mandibular bone resorption is present.The poor bone availability above the mandibular canal difficult the insertion of regular length implants 1,2 .There are different treatments for this clinical situation, including short dental implants (SDI), >6 to <10 mm in length, extra-short dental implants (ESDI), ≤6 mm in length 3 , or surgeries for vertical bone augmentation 2,4 .A recent systematic review showed at 1-year follow-up that SDIs have less morbidity, rehabilitation cost, and better survival rate (97%) than regular implants (92.6%) installed in a grafted bone area 5 .Besides, in this same study, the proportion of patients with biological and mechanical complications was lower for SDIs, with an incidence of 6%, while 39% of complications were reported for regular implants in grafted areas 5 .Meanwhile, over a 5-year follow-up period, it was shown that there was no statistically significant difference in implant survival and success rates between SDIs and regular implants in the grafted area 4 .Also, ESDIs compared to regular implants have similar survival rates, 96.2%, and 99%, respectively, as well as the technical complications incidence, 14.14%, and 18.36%, respectively, after 3-years of follow-up 6 .
In addition, a study that evaluated the long-term effectiveness of ESDI reported a survival rate of 94.1% at a five-year follow-up 1 .This slightly lower survival rate, when compared to regular implants, can be explained by its unfavorable biomechanics 7 , due to an increased crown-to-implant ratio (C:I) that creates a more significant vertical lever arm and a disadvantageous stress distribution 2 .These implants have a smaller bone/ implant contact surface, which leads to increased stresses at the bone and prosthetic components 8 .Therefore, SDI and ESDI generally have a wide diameter (WD) compensating the limitation in high, increasing the surface and its bulk, which improves the stress dissipation 9 , leading to better biomechanical behavior 10 .
The treatment plan also requires checking the patients' occlusion and the antagonistic type affecting the implant success 10 .In a physiological occlusion predominantly occur axial loads (AL), in the mandibular posterior region, transmitted by the long implant axis to the bone, resulting in an adequate stress dissipation 11,12 .However, when a non-physiological occlusion is present, the resultant occlusal force is an oblique load (OL), creating an unbalanced stress distribution 8 .Therefore, when the high C:I anchored by ESDI is used, the incidence of OL increases the bending moment of the vertical lever arm, causing a non-homogeneous force dissipation, leading to poor prognosis, which may contribute to peri-implant bone loss 8,12 .Clinical and in vitro studies showed that the increased C:I only negatively influences the stress distribution when an OL is present 8,13 .
Previous systematic reviews focused on C:I evaluation have shown no significant differences in biological complications and peri-implant health results 14,15 , being 2.36:1 the higher C:I evaluated 15 .Meanwhile, a recent four-year retrospective clinical trial concluded that the higher the C:I ratio (0.47 to 3.01), the less the marginal bone loss 16 .However, the biomechanical behavior of a challenging scenario where a 3:1 C:I crown supported by an ESDI, with 5 mm in length, at the severe reabsorbed mandibular posterior region, in the presence of OL, has not yet been investigated.That is critical since it can make the long-term success of this type of rehabilitation uncertain.
Besides, the benefits of using WD in ESDI have not reached a consensus in the literature since clinical and laboratory studies have not found differences in survival rates when assessing different diameters and lengths 2,17 .This fact contradicts the prerogative of better biomechanics due to its larger contact surface 10 .Therefore, there is a need for further studies to evaluate the rehabilitation mechanical behavior 12 before future prospective clinical studies.Thus, by using finite element analysis (FEA), the present study evaluated the influence of WD on the stress distribution of ESDI as support for single implant-supported crowns in the posterior region of the atrophic mandible, with a 3:1 C:I ratio, under AL or OL.For then, to verify if the WD is relevant enough to justify the insertion of an implant that will wear out more bone.The tested null hypothesis stated that WD would have no difference from the RD regarding the stress distribution.
Materials and Methods
Through the computer-aided design (CAD) software (SolidWorks; Dassault-Systemes SolidWorks Corp; Waltham, Massachusetts, USA) were created the 3D virtual models of a single crown, cement layer, cortical and cancellous bone.Also, CAD models of a universal abutment (4.5 x 2 x 6 mm) and two morse-taper implants of 4 x 5 mm (28.274 mm 3 , bone/implant contact surface: 101.39 mm 2 ) and 6 x 5 mm (75.75 mm 3 , bone/implant contact surface: 155.36 mm 2 ) were assessed virtually, and were left 2 mm submerged into the bone, which were obtained by the manufacturer (S.I.N Implant System, São Paulo, SP, Brazil).Two study factors were evaluated: I) implant diameter: 4 mm (RD: regular diameter, being this the control group) or 6 mm (WD) (Fig. 1); II) load angulation: AL or OL (30° off-axis) being applied at the mesiobuccal cusp (Fig. 2) 18 .The bone model had a 12.94 mm height and 16.11 mm of thickness, and a 2 mm layer of cortical bone surrounding the cancellous bone (Fig. 1) 19 .The crown of a mandibular right first molar, 13 mm in height with a 3:1 C:I 15 (Fig. 1), was virtually cemented on the abutment (70-μm layer), and four groups were created: RDAL (regular diameter implant under AL); WDAL (wide diameter implant under AL); RDOL (regular diameter implant under OL); WDOL (wide diameter implant under OL).The FEA models validation 20 was performed by past literature for the location of force application and bone layers dimensions and by past in vivo study for crown and C:I.
Braz J Oral Sci.2023;22:e238152 After assembly, the virtual models were exported to finite element software (ANSYS Workbench 15.0; ANSYS Inc; Canonsburg, Pennsylvania, USA) for a mathematical solution.A tetrahedral mesh was generated with an element size of 0.6 mm after convergence analysis with 5% tolerance.The Young modulus (GPa) and Poisson ratio (δ) of each material were set in the software according to table 1.The number of elements and nodes of each element is described on table 2. All components were considered homogenous, isotropic, and linearly elastic.Also, the contact conditions between implant/abutment were assumed as no separation, and the contacts crown/ abutment and implant/bone were assumed as bonded.Then, the models were fixed in two lateral portions of the bone segment and were submitted to a 200N load on the occlusal surface of the mandibular right first molar (Fig. 2) 8 .The equivalent von Mises stress (σ vM ) was used for the implant and the abutment 8,10 .Minimum principal stress (σ min ), and maximum shear stress (τ max ) 8,26 , were used for both cortical and cancellous bone.A qualitative analysis was performed for the implants, abutment, and bone, using the colors of the resulting FEA images.The colors varied from warmer (red) to cooler (blue) tones, with the peak stress represented by the warmest tone.
Results
Results for the FEA assessment are presented in table 3. Regardless of diameter, there was a significant increase in stress in all components, over 200%, under OL compared to AL results.Also, the stress was greater on the abutment and cortical bone and less on the cancellous bone and implant for WD groups.A substantial increase in stress was observed in cortical bone for WD groups compared to RD groups, being higher 66.3% for σ min and 99.8% for τ max under AL and higher 125.7% for σ min and 201.7% for τ max under OL (Table 3).For the AL groups, the peak stress concentration was in the area in contact with the apical region of the implant, being the maximum values found at σ min of 72.34 MPa (WDAL) (Fig. 3) and τ max of 42.02 MPa (WDAL) (Fig. 4).Meanwhile, in the OL groups, the highest stress concentration was in the cervical third of the bone, and the maximum values were at σ min 266.7 MPa (WDOL) (Fig. 3) and τ max 130.88 MPa (WDOL) (Fig. 4).
The analysis of σ min and τ max showed decrease stress in the cancellous bone for WD groups, about 44.9% for σ min and 55.9% for τ max under AL and 73.2% for σ min and 71.9% for τ max under OL (Table 3).Also, the images showed a peak stress concentration in the cervical third of the bone in all groups, and the minimum value of the σ min was 9.79 MPa (WDAL) and of the τ max 7.32 MPa (WDAL) (Fig. 5 and Fig. 6).Besides, the σVm evaluation images showed that in all groups, the peak stress area was at the abutment collar level (Fig. 7) and in the corresponding region of the implant (Fig. 8).The analysis demonstrated that with the WD, a low increase occurred in the abutment stress of 3.6% under AL (WDAL: 202.94 MPa) and 12.7% under OL (WDOL: 1157.4MPa) (Table 3).However, a decrease in the implant of 38.7% was observed under AL (WDAL: 185.98 MPa) and 38.2% under OL (WDOL: 873 MPa) (Table 3).
Table 3. Von-Mises criteria (MPa) for implants and abutment, minimum principal stress and shear stress for cortical and cancellous bone (MPa), and the differences between the groups and direction of the load.
Discussion
There is no consensus in the literature about the benefits of using WD in ESDI in the treatment of severe mandibular bone resorption in the posterior region 12 .Also, recent studies showed that a high C:I ratio only increases the stress concentration when OL is present 8,27 , being traumatic occlusion the primary cause of biomechanical complications 8,13,27 .Thus, by FEA, the present study evaluated the influence of WD on ESDIs stress distribution as support for single implant-supported crowns in the posterior region of the atrophic mandible, under AL and OL.The hypothesis that WD would have no difference from the RD regarding the stress distribution, had to be rejected.It was observed that WD in ESDI, under both load directions, showed a decrease of stress at the implant and the cancellous bone (WDAL: τmax=7.324MPa, σmin=9.795;WDOL: τmax=20.66MPa, σmin=25.23 MPa), a relevant increase in the cortical bone, and a possible slight increase in the abutment.Besides, when submitted to OL, there was an increase in stress in all components and groups by more than 200%, corroborating with previous studies 8,13,27 .
In this study, the stress distribution on the peri-implant bone was different when a WD was used.A relevant increase (up to 66%) in the stress can be observed in the cortical bone when τ max and σ min were evaluated independently of load angulation.This is important since some studies have reported, without a consensus, a critical threshold of compressive (ranging from 50 MPa to 170 MPa) and tensile stress (ranging from 34.72 MPa to 100 MPa) of the bone [28][29][30][31] , and in the WDAL, RDOL, and Braz J Oral Sci.2023;22:e238152 WDOL these values were overtaken.What shows the need for more studies and other methods of evaluation of bone impact when WD is used.Also, the figures in WD groups shows a stress peak in the cervical third of the bone of at least 311.5% under OL higher than the findings of the AL groups, which could be explained by the use of the WD implant providing a 34.73% higher bone/implant contact and wear on the cortical bone.These results corroborate with Elias et al. 27 , which evaluated the influence of the prosthetic crown height in SDI and found a higher stress concentration in the OL groups.
Meanwhile, in the WD groups, a decrease in the stress was observed in the cancellous bone, bringing the MPa values found within the limits of compressive and tensile stress at WDOL [28][29][30][31] .This may be related to its Young modulus, since its value is lower than that of the cortical bone.The greater the Young modulus the stiffer the material, the greater the stress accumulation 10 , and more resistance to deformation 32 .In the present study, when the WD implant was evaluated the contact between the implant and cortical bone was increased, leading to higher stress on the cortical bone and a reduction on the cancellous bone, which can explain the results 10 .This enhanced contact with the cortical bone may negatively influence the bone remodeling around the implants since the cortical bone is less vascularized than the cancellous bone, which leads to interference of blood supply that directly affects the bone resorption response 33 .According to the results of this study, this would only be a problem in the presence of oblique load.Considering that in the posterior region the pattern of forces is axial, perhaps it would not be a clinical problem, as long as the patient has a favorable occlusal pattern.
The consequences of higher stress concentration on the cortical bone associated with its decrease on the cancellous bone remain uncertain since low-stress values around the implant resulting in a bone loss due to disuse atrophy, while high-stress cause microfracture at the bone resulting either in bone loss or fatigue failure of the implant 32,33 .Also, since WD in ESDI increases the stress at the implant/cortical bone interface, being MPa values over the compressive and tensile limits of the bone [28][29][30][31] , it represents a potential biological risk for marginal bone loss that might be even higher under OL.Besides, the mechanical loading conditions regulate the morphology of the bone 34 , and it is still unknown how much bone/implant contact is necessary for the success of ESDIs 27 .
The results of von Mises stress showed, in all groups, a higher stress concentration at the surface of the abutment collar level and at the implant platform where it touches the abutment collar.In both loads, the WD showed an increase of 12.7% in stress at the abutment and a reduction of at least 38.2% in the implant.Despite this percentage difference, the color pattern exhibits a great similarity in the stress distribution in general for the abutment, and under axial load for the implant.This substantial stress reduction at WD implants might be explained by its structure 62% bulkier than RD implants.Since the stress increased over 400% at implant and abutment at the OL groups, clinically, would increase the risk of the implant, and abutment failure once was exceeded the limits of tensile yield strength 0.2% (483 MPa) and ultimate tensile strength (550 MPa) of the titanium grade IV 35 .Suggesting that should be avoided the use of ESDI when it is impossible to eliminate OL during mandibular excursive movements, for example, in a parafunction scenario.
Another important point to be highlighted is that a WD implant might reduce the bone mechanical resistance, since the remaining bone around it is reduced when compared to a RD implant.There is a literature gap regarding the effects generated by an overload on the cortical bone, when a mandibular implant-retained crown is evaluated under different load directions.Also, the maximum stress values of FEA studies strongly depend on the size of the mesh used.So, even with this study results being encouraging, showing that the WD ESDI can be a reliable option as shown in the AL groups, it also shows the necessity to perform further studies on this behalf.
Clinically the masticatory forces are not acting in just one way, and it is impossible to isolate the force direction.So, it is essential to perform in silico studies, which allow the researcher to evaluate and study every direction of occlusal forces like was performed in this study.Besides, the present study is a numerical theoretical analysis, and its results should be validated with an in vitro study assessing implant failure mode in the same conditions of this study.In addition, other simulations could be performed to estimate possible statistical differences, for example, by using different prostheses, abutments, and materials with different elastic modulus since they could reach a different result because of its dampers chewing loads 10 .Finally, a reliable way to effectively assess the influence on the bone would be performing randomized controlled trials.These studies must include patients with severe bone atrophy in the posterior region of the mandible with different types of occlusal patterns and a minimum of 1 mm cortical bone wall to surround the implant.
Therefore, extra-short implants with wide diameter result in better biomechanical behavior for the implant, but the implications of a potential risk of overloading the cortical bone and bone loss over time, mainly under oblique load, should be investigated.
Figure 4 .Figure 5 .Figure 6 .Figure 7 .
Figure 4. Maximum shear stress peak concentration for cortical bone (MPa) for all groups.Blue to red color represents stress values from lower to higher, respectively.
Figure 8 .
Figure 8. Von-Mises stress peak concentration (MPa) in implant.Blue to red color represents stress values from lower to higher, respectively.
Table 1 .
Mechanical properties of materials.
Table 2 .
Numbers of nodes and elements of each component.
RD, regular diameter groups; WD, wide diameter groups. | 4,333 | 2023-06-27T00:00:00.000 | [
"Medicine",
"Engineering",
"Biology",
"Materials Science"
] |
Cryptoasset Competition and Market Concentration in the Presence of Network Effects
When network products and services become more valuable as their userbase grows (network effects), this tendency can become a major determinant of how they compete with each other in the market and how the market is structured. Network effects are traditionally linked to high market concentration, early-mover advantages, and entry barriers, and in the cryptoasset market they have been used as a valuation tool too. The recent resurgence of Bitcoin has been partly attributed to network effects too. We study the existence of network effects in six cryptoassets from their inception to obtain a high-level overview of the application of network effects in the cryptoasset market. We show that contrary to the usual implications of network effects, they do not serve to concentrate the cryptoasset market, nor do they accord any one cryptoasset a definitive competitive advantage, nor are they consistent enough to be reliable valuation tools. Therefore, while network effects do occur in cryptoasset networks, they are not a defining feature of the cryptoasset market as a whole.
INTRODUCTION
The rapid appreciation and popularization of cryptoassets over the past few years has incited a large body of scholarship on understanding their behavior and their positioning in the market, particularly financial markets. As cryptoassets gradually became a household investment and transaction medium, they began to invite greater regulatory and investor scrutiny, which created the Permission to make digital or hard copies of part or all of this work for personal or classroom use is granted without fee provided that copies are not made or distributed for profit or commercial advantage and that copies bear this notice and the full citation on the first page. Copyrights for third-party components of this work must be honored. For all other uses, contact the owner/author(s). , © 2020 Copyright held by the owner/author(s). https://doi.org/10.1145/nnnnnnn.nnnnnnn need to better understand their function as a market of their own and as market that forms part of the greater economy. While early analyses focused on simple economic illustrations of the functioning of cryptoasset networks in isolation, later work started exploring market-wide phenomena, including the dominance patterns of some cryptoassets over others.
Since cryptoassets are based on blockchain networks and are therefore network markets, one important parameter that reflects and determines their behaviour is the relationship between their userbase and their value. This relationship has a long history in network markets under the theory of network effects. Network effects theory states that the value of a product or service is co-determined by its userbase . Then, for products or services that obey network effects, one can derive the value of the network, and therefore their relative value to each other too, for a given userbase assuming that the relationship between and is known, for example ∝ ( ), ∝ 2 , ∝ 2 etc. Initially, this insight attracted attention because of its predictive potential of cryptoasset valuation. Indeed a number of studies attempted to develop valuation models based on network effects that could be used by investors to predict the future value of their assets and the value of the market as a whole. However, the implications of network effects go far beyond valuation and, understood properly, they inform also the structure and competitiveness of the market making them a key input into policy-making and regulatory decisions. Most notably, markets that are characterized by network effects are commonly thought to be winner-take-all markets, where first mover advantage is key, entry barriers are high, networks hit tipping points of no return, and contestable monopolies or high concentration can be the natural state of the market. This is for two reasons: firstly, because the value of joining a network is increasing in the number of other network adopters, because the bigger the number of existing adopters the greater the utility every new adopter derives from it (pure network effects), and secondly, because for every new adopter joining the network, existing adopters also benefit (network externalities). In both cases bigger equals better (everything else equal), creating an incentive for users to join the network where the value will grow larger both for new and for existing users, which creates a snowball effect. This kind of power concentration in networks that exhibit network effects usually makes regulators uneasy, and therefore, if cryptoassets exhibit network effects, they would (and should) attract higher regulatory and investor scrutiny.
Extant literature on network effects in cryptoassets is limited and has focused almost exclusively on confirming or rejecting, usually for Bitcoin only, a specific application of network effects, namely Metcalfe's law, which states that the value of a network is proportional to the square of its users ( ∝ 2 ), and, if confirmed, it would be a useful valuation tool. However, this line of literature presents only a binary distinction between the existence or not of a specific type of network effects, focuses only on valuation, uses sub-optimal data, and has also been temporally limited to the period before the recent resurgence in mid 2019, or excludes periods, therefore missing key parts in the cryptoasset market evolution.
By contrast, our analysis takes a more comprehensive view of network effects in cryptoassets, and, while it confirms that network effects occur in cryptoassets, it shows that they do not have the usual implications associated with them in terms of according competitive advantages, resulting in market concentration, or serving as a reliable valuation tool. Firstly, we define network effects to occur when the value of a cryptoasset network changes supra-or infra-proportionately to changes in its userbase, thereby showing both positive and reverse network effects, while not being constrained by a specific version of network effects. We also use two proxies for value and userbase each to better capture what users perceive as the value of the network and how the network size (userbase) should be measured, and we base our results on cleaner vetted data. Moreover, we examine multiple cryptoassets to get a broader view of the industry, as opposed to previous works which focused on Bitcoin. Lastly, our analysis covers the entire history of the studied cryptoassets, which includes the valuation spikes and subsequent declines in 2014, 2017 and 2019. The spike in 2019 and the preceding decline from the heights of 2017 are particularly valuable because they help us show that the results obtained in previous studies which sampled only up to early 2018 do not hold based on more recent history.
BACKGROUND, MOTIVATION AND IMPLICATIONS
Network effects were first studied in the 1970s to more accurately capture the value and growth of telecommunications networks [29]. The intuition was that when the nature of a product or service is such that it relies on linking users together, the value of the product is co-determined by its userbase . More specifically, for every user added to the userbase of a product, value is created not just for the joining user but for existing users as well. As a result, each new user derives value from joining a network that is relative to the size of the network (pure network effects) and creates an externality in the form of value that is captured by the network of existing users (network externality). Conversely, for every exiting user, value is lost both for the exiting user and for existing users. This type of network effects was called direct network effects to distinguish it from later extensions to the theory, which accounted for the effects changes in the network's userbase have on complementary products and services developed for that network [8]. This latter type was called indirect network effects, and it is not the kind that will concern us here.
The powerful implication of (direct) network effects is the increasing returns to the userbase and ultimately to the product exhibiting network effects. Because for products that exhibit network effects every new adopter makes the product more valuable relative to existing size of its network, it creates incentives for other adopters to adopt the product with the bigger network over its competitors. Consequently, the more the userbase grows the more it invites further growth rendering the product increasingly more valuable and competitive. The exact relationship between value and userbase can vary; While one can say that in the most basic version of network effects the value of a product grows linearly with the number of users added to its userbase ( ∝ ) [33], most commonly network effects are used to describe relationships that are logarithmic ( ∝ ( )) [5], quadratic ( ∝ 2 ) [24] or other (e.g. ∝ 2 ).
Network effects have found application in numerous industries and business models ranging from telecommunications [4,16], to web servers, PC software [17], airline reservation systems, ATMs [14], and platform systems [7]. Indeed, the intuition and implications of network effects have been so pervasive that they have been invoked in any industry where the consumption or use of a product by consumers makes the product more valuable for others (for a collection of relevant literature see [19]. It is no surprise that cryptoassets have also been hypothesized to exhibit network effects. The combination of the inherent network nature, the meteoric rise in popularity (read: userbase), and the substantial price volatility (read: value) has suggested a strong-if elusive-relationship.
The particular motivation behind the study of network effects in cryptoassets has so far been to discover a valuation formula: if we know the function between userbase and value, then with informed guesses on the network's growth we can predict future prices [27,31,34]. But valuation formulas reduce network effects down to a binary distinction represented by a single function. While useful as prediction tools and high-level descriptors of cryptoasset trends, valuation formulas provide little granularity.
Our motivation and goal is, instead, to provide more high-level view of how network effects influence the cryptoasset market as a whole, and particularly what they say about the potential for concentration in the market and about competitive (dis)advantages of one cryptoasset over others. These are the most impactful implications of network effects, and they are desirable for those networks that can exploit them, but undesirable for their competitors or for regulators who have to deal with concentrated markets. We work with numerous cryptoassets so that we can obtain a market-wide overview (limited by how big and representative our sample is), and we study them from their inception until early 2020 which allows us to capture all historically important phases, including the resurgence in 2019, which extant literature has not had a chance to consider. This type of approach allows us to draw insights about the structure and competitive dynamics of the cryptoasset market. It goes back to the early wave of "Bitcoin maximalism", which stood for the idea that the optimal number of currencies as alternatives to the mainstream financial system is one, and altcoins will eventually be rendered obsolete as more and more users gravitate toward the biggest, most stable, most widely accepted cryptocurrency, namely Bitcoin. At the time, Bitcoin maximalism was rejected by Vitalik Buterin, the creator of Ethereum, correctly pointing out that the cryptoasset universe is not a homogeneous thing, and that therefore there is no one single "network" around which network effects would form [6]. We expand on that thinking.
Looking at network effects to study the competitive dynamics of the cryptoasset market and its potential to concentrate around one or a small number of cryptoassets can provide useful insights for industrial policy. Normally, a showing that cryptoassets exhibit network effects would suggest that early cryptoassets have a firstmover advantage and may lock the market in [13,18,23], even if they are intrinsically inferior to other comparable cryptoassets [5,15,20]. While, the market seems to have moved away from that danger, network effects theory also suggests that, assuming homogeneity, once a cryptoasset hits a tipping point, it may fully prevail because new users will always prefer the cryptoasset with the larger userbase (the so called "winner-take-all" markets, which Bitcoin maximalism relied on) [13,23]. Homogeneity is, of course, a matter of degree, and it is still likely that, if a cryptoasset exhibits stronger network effects than its peers, it can prevail at least within a sub-segment of the market. The flip side of network effects can also be observed, whereby the loss of a user results in a supra-proportionate loss of value (i.e. more value than the user intrinsically contributed individually), which incites further losses and so on. This means that rapid depreciation is more likely in cryptoassets characterized by network effects. The rapid appreciation and depreciation cycles coupled with the winner-take-all characteristic can in turn result in cryptoasset markets that are successively dominated by a new winner in every era (successive contestable monopolies). Then, if this is the natural state of the market, artificially forcing more competition may not be optimal.
These insights are well-applicable in financial markets. For instance, the influential "Cruickshank report", an independent report on banking services in the United Kingdom prepared for the UK Treasury, which has in turn influenced regulatory and legal decisions [1, 2], warned about the far reaching implications of network effects: "Network effects also have profound implications for competition, efficiency and innovation in markets where they arise. Establishing critical mass is the first hurdle, as the benefits to customers and businesses of a network arise only gradually with increasing use. It is possible to imagine a world in which electronic cash is widely held and used, for example, but much harder to see how to get there. Once a network is well established, it can be extremely difficult to create a new network in direct competition. ... Where network effects are strong, the number of competing networks is likely to be small and the entry barriers facing new networks will be high" [11]. As the fintech industry is heating up, network effects have also been cited there as a strong factor in entrenching existing market power of financial services (see e.g. the recent proposed acquisition of Plaid by Visa [12]), and such risks have also been highlighted in the cryptoasset market, with models showing that certain conditions can allow cryptoasset markets to become oligopolies and market players entrench their position in the market [3,10].
PRIOR LITERATURE AND CONTRIBUTION
A number of papers have investigated aspects of the application of network effects in cryptoasset networks. The focus has been to determine whether the value of cryptoassets (and mainly Bitcoin) complies with network effects, and in particular on whether it follows Metcalfe's law, which is the most popular iteration of network effects and stipulates that the value of a network grows at a rate proportional to the square of the number of users ( ∝ 2 ).
The early influential analysis by Peterson [27] remains the point of reference. Peterson developed a valuation model for Bitcoin's price based on Metcalfe's law for the period 2009-2017, using wallets as a proxy for users, Bitcoin prices as the proxy for value, and a Gompertz function to account for growth. He found that the price of Bitcoin follows Metcalfe's law with R-square of 85 percent. In a revised version of the original paper that extends through 2019, Peterson re-confirms the application of Metcalfe's law to Bitcoin [28]. However, he excludes significant periods of time on the grounds of price manipulation, during which the value of the Bitcoin network, as measured by the USD price of Bitcoin, lies well outside of Peterson's model predictions. Van Vliet [34] enhanced Peterson's model by incorporating Rogers' diffusion of innovation models to better capture population parameters and growth rates. By doing so, van Vliet raised R-squared to 99 percent. Shanaev et al. [31] acknowledge the utility of Peterson's and van Vliet's analyses but depart from them in that their model does not rely on historical data for the estimation of the coefficient of proportionality, which raises an endogeneity problem. They still use Metcalfe's law but only as only as one of the building blocks of their model. Civitarese [9] rejects the applicability of Metcalfe's law to the value of the Bitcoin network by running a cointegration test between price and an adjusted number of wallets' connections.
Gandal and Halaburda [18] use a completely different approach to examine the existence of network effects in cryptoasset networks. They define network effects as the reinforcement effects the price of a cryptoasset has on the price of another cryptoasset. With Bitcoin as the base cryptoasset, the idea is that, if network effects are in place, as Bitcoin becomes more popular (price increase), more people will believe that it will win the winner-take-all race against other cryptoassets resulting in further demand and higher prices. Therefore, network effects would manifest themselves as an inverse (negative) correlation between the prices of the sampled cryptoassets. For the period May 2013 -July 2014, their results showed signs of network effects after April 2014.
Our analysis complements and differs from prior literature in several ways. Firstly, we do not focus on a specific network effects formula; we rather look at when, to what degree, in which cryptoassets, and for what proxies of value and userbase network effects are observable (defined as supra-proportional change in value relative to userbase) regardless of which particular curve/function they follow. Secondly, we go beyond Bitcoin to examine six cryptoassets that we have selected as representative of different features and characteristics to better be able to observe potential industry-wide trends. This helps us notice whether one cryptoasset has the potential to dominate the market or multiple cryptoassets benefit from the same network effect forces. Thirdly, we use different parameters as proxies for value and userbase to more fully capture the functionality and usage of cryptoassets in the market. Importantly, we do not rely on the total number of users as a proxy for userbase like extant literature, because many of those addresses are dormant or permanently inaccessible and therefore economically irrelevant. Fourthly, we study the full history of cryptoassets from their inception to today which allows us to observe their different phases, including the price collapse in 2018 and the resurgence in mid-2019, which dramatically change the picture of network effects and which have been missed by previous studies. Lastly, we work with data sets that have been meticulously cleaned to filter out spurious or manipulative activity, which improves the accuracy of our results compared to data-sets that are pulled unfiltered from the network. Our analysis confirms the existence of network effects, but also that they do not have the results usually associated with them on the market.
METHODOLOGY AND DEVELOPMENT
We study the application of network effects in Bitcoin (BTC), Dogecoin (DOGE), Ethereum (ETH), Litecoin (LTC), XRP and Tezos (XTZ). The selection of these cryptoassets was made on the basis of diversity and feasibility. We aimed to study cryptoassets that exhibited different attributes in terms of age, market capitalization and any special features that make them stand out from other competing cryptoassets in order to build a representative sample of the crypto-economy [22]. We also limited the study to cryptoassets for which we could get reliable, standardized time-series data from the cryptoassets' initial release to the time of the study [25]. The unreliability of the prices reported by exchanges in the early days of the industry led us to consider Bitcoin from July 2010, Litecoin from March 2013, and XRP from August 2014-the rest from their beginning. Table 1 summarizes the attributes of each chosen cryptoasset. We first define network effects. Network effects occur where the value of the network grows supra-proportionately to the number of users that participate in the network. Reverse network effects occur where the value drops supra-proportionately to the number of users that leave the network. Unless there is a reason to distinguish between positive and reverse network effects, we collectively refer to them as network effects. Therefore, we define network effects to occur in cryptoassets when a positive value change Δ > 0 is larger than a positive userbase change Δ > 0, or when a negative value change Δ < 0 is smaller than a negative userbase change Δ < 0. Notice that we do not consider that network effects apply when value and userbase move in different directions, e.g. when the value increases while the userbase decreases, regardless of which increases or decreases more.
Thus, network effects occur if In our analysis we define change at time similar to log returns, i.e.
+1
(1) (2) Then, we identify appropriate proxies to represent value and userbase . To represent we use two proxies: (a) token price and (b) transaction value. The two proxies represent different aspects of the value users assign to cryptoassets. In theory, even one proxy applied to one cryptoasset would be enough to demonstrate (or not) network effects (as has, for example, been done in previous literature that relied only on token price), assuming the proxy and cryptoasset are representative. However, because cryptoassets are differentiated resulting in diversified usage patterns, and because the chosen proxies express different ways by which users perceive the value of the network, a multitude of cryptoassets and proxies was used in an effort to better represent the industry.
Token Price (PriceUSD): The first parameter we use is token price, which is the fixed closing price of the asset as of 00:00 UTC the following day (i.e., midnight UTC of the current day) denominated in USD (for a detailed explanation of Coin Metric's methodology on toke price see [25]). Token price expresses value in terms of market forces, namely the point at which supply meets demand. It is the value that users as market participants collectively assign to a given cryptoasset by deciding to buy and sell at that price level. We assume that the studied cryptoassets trade under normal market conditions; any acknowledgement of price manipulation that may have occurred at times has been accounted for in the cleaning of data by Coin Metrics [25].
Transaction Value (TxTfrValAdjUSD): The second proxy of choice is transaction value, which expresses the USD value of the sum of native units transferred between distinct addresses per day removing noise and certain artifacts to better reflect the real economically relevant value circulating in the network. The assumption is that as the network becomes more valuable to users, they will use it more frequently and/or to transfer greater value among them. Therefore, transaction value as a proxy sees cryptoassets as means of transaction. We considered and rejected transaction count as an appropriate proxy, because on some networks a large number of recorded transactions are unrelated to value transfer, but rather to the operation of the network, e.g. consensus formation on Tezos [26]. One could retort that even these non-value-carrying transactions reflect engagement with the network and that therefore are an indication of the value of the network to users. Even so, lumping together value-carrying and operational transactions would taint the comparison across cryptoassets, since on some cryptoassets the majority of transactions are operational (e.g. Tezos, see [26]), while on others value-carrying (e.g. Bitcoin).
Next, to represent we select the following proxies: (a) addresses with non-zero balance (b) trailing 6-month active addresses and . Different ways to represent userbase more fully captures the relationship between value and userbase. We considered and rejected counting userbase based on total number of addresses (like all previous literature), because of the large number of inactive addresses.
Contrary to other industries where network effects have been studied and where inactive users are eventually purged from the network (e.g. mobile phone subscriptions, social networks), so that total user count may still be a good approximation of the economically meaningful userbase, this is not the case with cryptoassets. Instead we opted for two variants of addresses with non-zero balance, as defined below.
Addresses with Non-Zero Balance (AdrBalCnt): This proxy represents the sum count of unique addresses holding any amount of native units as of the end of that day. Only native units are considered (e.g., a 0 ETH balance address with ERC-20 tokens would not be considered). The utility of this proxy lies in that it excludes all noneconomically active addresses, the assumption being that addresses with zero balance are dormant (similar to bank accounts with zero balance). This choice responds to criticism that has been raised with regard to extant literature that tended to use all addresses or wallets as a proxy for users. Despite our choice of improved metric, it still remains a fact that there is no one-to-one mapping between addresses and actual users, which is a common problem to any network or service, e.g. the same person may have multiple bank accounts. While there are methods to de-cluster actual users from wallets and addresses, these are not sufficiently precise and are unavailable or inapplicable across cryptoassets [21]. We also acknowledge that on networks with lower transaction fees it is easier to generate and/or maintain addresses with balance, and to counter that we could raise the amount of native units the counted addresses should have, but this would introduce a subjectivity question without even fully eradicating the initial problem of spurious addresses.
Trailing 6-Month Active Addresses (6MAdrActCnt): This proxy counts all unique addresses that have been active at least once over the trailing 6-month period from the time of measurement. Repeat activity is not double-counted. Traditionally, most userbase measurements are taken in time frames that range from one month to one year. Given that cryptoassets are of relatively young age, which may suggest that their userbase is expected to interact with them less frequently, and that part of their utility involves simply owning them, which does not generate any activity, we decided that a 6-month time frame sufficiently captures active userbase.
Before we derive network effects, we first calculate the Pearson correlation between value and users which is informative in terms of their overall relationship. Next, we obtain relevant measurements of network effects. We rely predominantly on the PriceUSD-AdrBalCnt pair of proxies for value and userbase, but additional measurements are in the Appendix. To see how prevalent network effects are in the studied cryptoassets we calculate the ratio of total days to the days where network effects were observed (separately for positive and reverse) for each cryptoasset. To see how strong network effects are we calculate the ratio of total days to the sum of the network effects observations over the days they occurred for each cryptoasset (separately for positive and reverse). To see how strong network effects are in cryptoassets relative to each other we reduce to a 100 day period. The results are presented in Part 5 and the analysis of the results in Part 6. Metric meaning PriceUSD Token price TxTfrValAdjUSD Transaction value AdrBalCnt Addresses with non-zero balance 6MAdrActCnt
Metric abbr
Trailing 6-month active addresses NFX Network effects Table 2: Legend of metrics in use.
RESULTS
We are looking for network effects in the relationship between value and users of various cryptoassets as represented by the proxies defined previously. Four pairs (2x2 proxies) are possible: • Token Price -Addresses with Non-Zero Balance: This pair demonstrates network effects expressed as the change of monetary value of a cryptoasset relative to the users that hold any amount of that cryptoasset. By counting only accounts with non-zero balance, we filter out economically dormant users. • Token Price -Trailing 6-month Active Addresses: This pair demonstrates network effects expressed as the change of monetary value of a cryptoasset relative to the users that have been active at least once in the trailing 6month period on that cryptoasset's network. Counting all active users over a recent time segment (usually 1, 6 or 12 months) is a common measurement of network or platform userbase and less conservative than daily active users.
• Transaction Value -Addresses with Non-Zero Balance:
This pair demonstrates network effects expressed as the change of transaction value of a cryptoasset relative to the users that hold any amount of that cryptoasset. • Transaction Value -Trailing 6-month Active Addresses: This pair demonstrates network effects expressed as the change of transaction value of a cryptoasset relative to the users that have been active at least once in the trailing 6month period on that cryptoasset's network.
Before we derive network effects, we calculate, based on the above pairs, the Pearson correlation between value and users which tells us whether, as a general matter, cryptoasset value and userbase are moving in the same direction. This already provides an indication of whether cryptoassets become more valuable as their adoption increases.
It is evident that only BTC shows a strong correlation between value and userbase, at least when userbase is measured by our main proxy of total addresses with non-zero balance (AdrBalCnt), with LTC showing the next highest correlation, which is, however, average and only holds when value is measured as value in fiat currency (PriceUSD). Correlations when userbase is measured as addresses that have been active in the trailing 6-month period (6MAdrActCnt) tend to be higher although still not consistently so. Higher correlation using 6MAdrActCnt might be explained on the grounds that user activity picks up during phases of large price movements. Overall, the mediocre and inconsistent correlations between value and userbase provide a first indication that a blanket Table 3: Pearson correlation between value and user proxies conclusion that the cryptoasset market is characterized or not by network effects is unwarranted. Next, we obtain relevant measurements based on the PriceUSD-AdrBalCnt pair of proxies for value and userbase as presented in Table 4 (additional measurements for other pairs are in the Appendix). As explained in the methodology, we believe these are the most appropriate proxies. Column 5 of Table 4 shows prevalence of network effects for each cryptoasset as calculated by the ratio of total days to the days where network effects were observed (separately for positive and reverse). Column 6 of Table 4 shows strength of network effects as calculated by the ratio of total days to the sum of the network effects observations over the days they occurred for each cryptoasset (separately for positive and reverse). Column 7 of Table 4 shows relative strength of network effects across cryptoassets by reducing to a 100 day period. This allows us to compare how strong network effects are across cryptoassets regardless of how prevalent they are across them.
ANALYSIS
Our results are useful in reaching a number of conclusions on how network effects inform the structure and evolution of the cryptoasset market.
(1) Network effects do not provide precise valuation predictions: The most common application of network effects theory has been to draw insights into future cryptoasset pricing based on the evolution of their userbase. Our results indicate that network effect observations in cryptoassets are frequent but inconsistent and therefore they cannot be relied on, generally, as a valuation tool as previous literature suggests (Figures 2 and 3). They are most frequent in XRP (45 percent of time in the pair Token Price-Addresses with Non-Zero Balance) and least frequent in LTC (29 percent of time in the same pair). While they appear more consistent in ETH and XRP, their results can be somewhat misleading at first glance: ETH's and XRP's userbase (AdrBalCnt) was constantly increasing and so any supra-proportionate increase in price registered as a (positive) network effect observation (blue lines in (c) and (e) in Figure 2). However, the positive network effect observations are frequently punctuated by days/periods of no network effect observations during which the price either does not rise supra-proportionately to userbase or drops. In cryptoassets such as BTC and LTC, where userbase fluctuates, it is easier to notice the changes in network effects trends (blue and red lines in (a) and (d) in Figures 2 and 3), even through network effect frequency is comparable to ETH and XRP. Therefore, it is hard to conclude that in any cryptoasset network effects exhibit constant patterns that, if extended into the future, can hold predictive value. This does not mean that we do not acknowledge the exponential long-term price increase of some cryptoassets (Figure 1), but we note that this is not linked consistently to their userbase growth, which is what network effects theory suggests. [28]. A third explanation relates to the proxies used. Some previous studies rely on wallets (total addresses) as the proxy for userbase, which is a more crude measurement than our preferred addresses with non-zero balance, as the latter show only economically active users and are therefore a better approximation of relevant userbase.
(2) Reverse network effects are also noticeable meaning that cryptoassets are vulnerable to rapid decline, not just conducive to rapid growth: While network effects have mostly been used to describe growth patterns, they are equally applicable in describing decline. Reverse network effects reflect situations where a decrease in users is linked to a larger decrease in value. Such observations are important, because they show that each user loss incurs a greater loss of value and therefore expose the potential for a rapid decline of the network once user exodus begins. Reverse network effects therefore highlight the precariousness of success (as measured by proxies of value). Most cryptoassets exhibited at least one prolonged period where reverse network effects were dominant, during which phases their value contracted disproportionately to the contraction of their userbase ending up less valuable than their userbase size would otherwise suggest or mandate during that period. This is noticeable both when userbase is measured by addresses with non zero balance, but it is even more pronounced when userbase is measured as trailing 6-month active addresses (Figure 3). This makes sense since the users active in the trailing 6-month period are more likely to be responsive to price fluctuations compared to users who simply hold some balance on their account. From Figure 3 it is also evident that user disengagement is almost consistently observed after every price crash (as manifested through the reverse network effects that begin 6 months after many of the crashes), and the fact that price continues to decrease supra-proportionately to userbase, as measured by active users in the trailing 6-month period, 6 months after the crash, may be indicative of the lasting effects user exodus has on the value of cryptoasset networks. Generally, however, while reverse network effects serve as a cautionary note that rapid decline of value can be triggered by user exit, they are weaker in magnitude than positive network effects (Table 4). So, overall, positive network effects (albeit inconsistent) still seem to characterize cryptoasset networks.
(3) Cryptoassets do not seem to be a winner-take-all market: A common corollary of network effects is that they eventually cause the market to gravitate toward oligopolistic structure, since, everything else equal, users prefer to join the network where the value from their joining will be maximized. This causes a "rich-get-richer" effect where the most valuable network continues to become even more valuable as users prefer to join that over others. Such markets tend to become oligopolistic, with the usual downsides of such industry structure (higher prices, reduced output, entry barriers; lower variety and innovation), and can therefore be a cause for concern. For this to be more likely to happen the various networks (=cryptoassets) must be undifferentiated and switching among and multi-homing across networks must be rare or costly [30]. These features do not seem to characterize the cryptoasset market, which accordingly appears less susceptible to a winner-take-all trend, at least on account of network effects. Indeed, of the thousands of available cryptoassets many serve different purposes, and users can own multiple cryptoassets at the same time and enter and exit their networks without friction. As evidenced by our results, the fact that the various cryptoassets we studied exhibit network effects of comparable relative strength (Column 7 in Table 4), and that they retain their userbase and valuation cycles ( Figure 1) seems to suggest that the underlying market features, including network effects, do not lead it toward an oligopolistic structure.
(4) Network effects strength across cryptoassets is comparable and therefore network effects do not accord a single cryptoasset a strong comparative advantage over its peers, undermining fears of concentration: Besides frequency and duration, i.e. what period of a cryptoasset's lifetime is dominated by network effects, another useful parameter of network effect observations in cryptoassets is their strength, i.e. the magnitude of the impact of a userbase change to value change [32]. Strong network effects can be indicative of higher homogeneity or cohesion within the network, where the addition of each new user (e.g. investor) affects existing users of that closely-knit network more than if it was a different looser network. In turn, this is reflected in the value of the network, or they may be indicative of stronger reputational effects, where the addition of each new user signals major changes for the network, which are then reflected in its value. Our results show that the comparative strength of network effects across the studied cryptoassets is similar (Table 4). This leads us to believe that no single cryptoasset benefits from network effects significantly more than its peers and therefore that no cryptoasset enjoys an overwhelming competitive advantage over its peers on account of network effects. A necessary corollary observation is that network effects accrue at similar levels to the studied cryptoassets, which means that network effects as a phenomenon, characterizes the cryptoasset industry as a whole (at least based on our sample), not just Bitcoin, which has been the main subject of many of extant studies in the area. This is not a surprising finding, but it is worth highlighting that it lends support to the previous point that the structure of the cryptoasset market does not seem to be such where network effects lead it to concentration around a small number of cryptoassets or that it helps cryptoassets overtake their peers on account of network effects. This is most likely because cryptoassets are differentiated and multi-homing and switching are pervasive.
(5) Network effects are not consistently observed during the early days of cryptoassets and therefore it is doubtful that they can be relied on as a tool to bootstrap a new cryptoasset: A common business model when launching new products or services in digital markets is to exploit network effects to quickly establish a growing foothold. Particularly if the product or service is also the first of its kind to hit the market, network effects can dramatically augment the first mover advantage, everything else equal. Our results indicate that network effects are not consistently observed in the studied cryptoassets during their early days (the first year of data); in particular, DOGE, XTZ and LTC do not exhibit consistent positive network effects neither by token price (PriceUSD) nor by transaction value (TxTfrValAdjUSD) as proxies for value (Figures 2 and 5). The lack of consistency is even more pronounced when userbase is measured by active addresses in the trailing 6-month period, which is an instructive measure here, because it tracks recent user activity which is the driver of early adoption. In Figure 3 only BTC and ETH have a claim to positive early network effects and in ETH they are sparser. This suggests that new cryptoassets cannot necessarily hope that network effects will assist in their initial uptake. It is useful to dispel this hypothesis because investors are looking for patterns in events that may trigger valuation changes (e.g. the hypothesis that cryptoasset value as measured in monetary terms increases once the cryptoasset is listed on a major crypto-exchange).
(6) Comparison between network effects on price and transaction value reveals sensitivity to price, which can be a competitive disadvantage: Extant literature has relied exclusively on token price as the proxy for network value. Using transaction value too helps us draw useful comparisons. For this, it is most instructive to rely on trailing 6-month active addresses as the proxy for userbase, because this proxy is more responsive to value fluctuations. Then, a comparison between the strength of network effects measured by token price (PriceUSD) and by transaction value (TxTfrValAdjUSD) reveals that some cryptoassets experience greater fluctuations in their transaction value relative to their token price. During upturns, network effects tell us that token price and transaction value increase more than the userbase increases, and during downturns, reverse network effects show the opposite. By comparing the ratios among cryptoassets of the sum of network effects when value is measured by token price and the sum of network effects when value is measured by transaction value one can observe differences in how transaction value is affected among cryptoassets. Specifically, the ratios for BTC, DOGE, ETH and LTC are similar ranging from 0.12 to 0.14 for positive network effects and 0.07 to 0.09 for reverse network effects, whereas XRP's is 0.07 and for XTZ's is 0.06 for positive network effects and 0.04 and 0.03 for reverse network effects (compare sum ratios in Figure 3 and Figure ??). This means that during periods of positive network effects, XRP's and XTZ's transaction value grows more than their token price grows relative to their userbase, and that during periods of reverse network effects, XRP's and XTZ's transaction value drops more than their token price drops relative to their userbase. This kind of increased volatility may be generally undesirable, but it is particularly problematic during downturns (reverse network effects) because it shows that activity on XRP and XTZ networks is more drastically affected making them more sensitive and less resilient, which is a competitive disadvantage. Our results hold too when we look exclusively at 2017 and 2018 as the years with the most sustained price increase and decrease respectively.
CONCLUSION
Network effects can be among the most common and influential factors shaping market dynamics in industries where products and services are built around networks. It is no wonder that they have been cited as a determinant in how cryptoassets grow in value and compete. Our analysis show that while network effects do characterize cryptoassets, they do not result in the usual concentration and competitive advantage implications usually associated with them. Our work also invites further research to determine the exact scope and conditions under which network effects apply. More precise proxies for userbase and value and accounting for exogenous effects are steps in the right direction. | 9,816.8 | 2021-01-15T00:00:00.000 | [
"Computer Science",
"Economics"
] |
Research on the Aim Value ’ s Variation of Yen Fu ’ s Constitutionalism Political Thought
The aim value for the intellectuals’ constitutionalism political thought shared the variation attribute under the circumstance that it was no great change for three thousand years at the end of the Qing dynasty. It was appeared for the value of freedom and order in the early time of Yen Fu thought, but the freedom and order as the implicit aim value were covered by the wealth and power as the visible aim value in the condition that the national salvation dominated all else. In a word, there are a vertical static value genealogy and horizontal dynamic value relation in which the wealth and power, freedom and order are as the aim value and have been changed with the passing time in his constitutionalism political thought, that is to say, his ideal that surpasses to simply pursue the value of wealth and power is to never give up.
Introduction
Firstly, the question about the aim value for the intellectuals' constitutionalism political thought under the circumstance that it was no great changes for three thousand years was appeared.For the intellectuals who were lived in the end of the Qing dynasty, their first task was to seek the national wealth and power and national revival.So, the constitutionalism politics that was regarded as a cultivated thought to make country become richer and make people become stronger was brought into China.Did the intellectuals who supported the constitutional politics during that time only care about the wealth and power?Or, were the wealth and power, the ultimated value, in their minds?Later, many intellectuals advocated the civil rights and the regional autonomy, which were not just a tool or a means in their view.They would pay less attention to the value of wealth and power in a certain events or a period of time that the intellectuals were experienced.However, it is becoming more neces-sary to understand and learn the recent constitutional practices for one hundred years which have contained the Chinese classical legal tradition, the socialist legal tradition and modern Western legal tradition for the research of the other aim value on what is the other aim value and why we pay attention to them.Secondly, the selection for the research object.Why does the thesis regard the thought and practice of Yan's constitutional politics as the studied core?At the end of Qing dynasty which was visible for the internal disturbance and external aggression in China, his thought and practice about the constitutional politics were an outstanding example among them when the early modern intellectuals began to seek for the difficult road of the constitutional politics.More narrowly conceived, as professor Shi Huazi pointed out that he was standing out of political actions.His practice about constitutional politics was achieved by commenting rather than practicing.More broadly conceived, he himself has also taken part in a series of activities about constitutional politics, such as teaching the Western learning, assisting to build a Russian school in Tianjin, creating a daily-Guo Wen Newspaper and a weekly-Guo Wen Hui Bian, and translating masterpieces.However, it is still obvious to see for his negative actions and he has even been accused of Being a Talker rather than an Actor, so the research value of his thought is much higher than that of his practice.This paper will divide his thought and practice about the constitutional politics into three periods to study.The early stage is from 1895 to 1897; the mid stage is from 1898 to 1910; the late stage is from 1911 to 1921.
To Seek for the Wealth and Power in the Early Stage
Under the historical background that the national capitalism has developed, the national crises is gradually deepened and the government rule is becoming more and more incompetent at the end of Qing dynasty before the reform movement of 1898, Yen fu like other intellectuals of western and reformists all hoped to achieve the national prosperity through establishing the constitutional politics, revitalizing the business and resisting the foreign aggression.Through his accumulation of learning and experience for several years, Yen fu has published a series of fierce political essays such as On the Urgency of the Changing World in 1895 which expressed his initial thought about the constitutional politics-called freedom and democracy out and formed the thought about constitutional politics that treats the wealth and power as the aim value.
Firstly, he pointed out that the pattern difference to run the country between China and Western countries would lead to the national condition run the opposite direction of achieving the wealth and power.The great difference to run the country between China and foreign countries is that Chinese paid more attention to yesterday while Westerns laid more emphasis on today... (Wang, 1986) [1].On one hand, the Chinese lay stress on copying old things while Westerns focus on innovation to make the knowledge update in academic aspect.The Chinese accept the cycle fatalism-peace and riot or ups and downs are cycling in a country, yet Westerns work hard to explore the political system that will keep the country prosperous for a long time on the system of governance.On the other hand, the ruler in the old time selected governors through the imperial examination to lecture the people in culture and morality.And the imperial examination itself only emphasized the moral cultivation and the ruling means of candidates, which finally led to the disappearance of people's vitality and initiative and made the social life lie in a low level for a long time to maintain the low level of ruling order.Based on the contrast between these two governance pattern, he considered that the result could only be led the economic development remain stagnant for long term and their vitality is unable to be stimulated by the old system as the people is the passive role in old China.There is a mushroom growth in the western countries after the constitutional politics is established.
Secondly, the co-governance with people is realized to finally achieve the wealth and power only if the constitutional politics is established.Of course, there are a prerequisite for the co-governance with people.It would be the first task to give the people freedom under the circumstances of the people's intellectual level at the end of Qing dynasty, as Yen Fu once put that there are differences between freedom and without freedom.It is the basis for giving people freedom that make people become more intellectual, more strong and health and more moral.The basic concept of constitutional politics can be accepted by People to understand the relevant knowledge about constitutional politics through making people become more intellectual, and the social consensus can be achieved before the political reform to pool the strength of society.It is the way to solve the difficult situation that China was forced to be an opium importing country and the weakness of people and army led to the failure of war after two opium wars for making people more strong and health.It is for making people more moral to cultivate the citizen's subject consciousness and public awareness which the constitutional politics need.The soul of science is pursuing the truth and the core of democracy is not autocratic, which are similar to the Chinese culture, but why can't we carry out these thought like western countries (Wang, 1986) [1]?This is because people do not change from a subject to a citizen, they do not form the subject consciousness which includes rights, freedom, equality and participation and the public awareness which contains laws, negotiation, the social morality and patriotism.The second task is to perfectly deal with the relationship between the emperor and the people, as Mencius once said that the people are the most important element in a state, the next are the gods of land and grain and the least is the ruler himself.The congress is established to act as the bridge to communicate and co-govern between the rulers and the subjects.
Relying on Freedom in the Medium Term
He has chosen the translation as his lifelong career in the mid stage is from 1898 to 1910.A series of masterpieces translated by Yen Fu has become the classical literature for later constitutional politics.Freedom in this time has not been the way to realize the wealth and power, but has become the aim value of his thought about constitutional politics.Three important elements to realize the freedom are the freedom's condition, conduct and safeguard that are all presented in the famous books translated by Yen Fu and form a distinctive value system about freedom.
Firstly, the conditions to realize the freedom.Although there are no author's note in translations for John Mill's book-on liberty, he explained the conditions of freedom in his translation book in details that everyone has freedom after entering the society, but if there are no limitations for freedom, the society will be full of conflicts.So there must be a dividing line among everyone's freedom, which is the same as the moral standards in Great Learning that talented people rely on to make the society peaceful (Wang, 1986) [1].He put freedom in the society and professionally explained the conditions of freedom from a passive perspective and he thought that freedom is a space for people to avoid different unreasonable limitations and to choose by themselves.
Secondly, the realization of freedom needs the initiative action of people.The rights of freedom for citizens may be equal in law, but the value of freedom for them is different among them.Yen fu has said that the people fight for freedom from aristocrats under the rule of aristocrat and people fight for freedom from emperors under the rule of emperor, but under the system of constitutionalism and democracy, people will not fight for freedom from aristocrats and emperors who are limited in law at this time and they should strive for freedom mainly from the society, the nation and the current fashion (Wang, 1986) [1].He points out that no matter aristocrats, emperors or society, nation and current fashion all should stop outside the field of personal freedom.
Thirdly, the realization of freedom needs the protection of political system.The political freedom is established on the certain separation of legislative power, administrative power and judicial power (de Montesquieu, 1959) [2].According to the theory of dividing power proposed by Montesquieu, Sun Yat-sen came up with the separation of power and divided the state political power into two parts.The one is Quan or called regime, the other is Neng or called governance.Although the personal freedom is endowed by god, only by reasonably and effectively controlling the power of government can the individual freedom finally realize.
Expecting the Order in His Late Period
When Yen fu came into his old age, his political thought about constitutionalism was also in the late period.He, at this time, as professor Zhang Kaiyuan described that constitutionalists and the upper merchants of southeast area always lack the courage of living on their own, they reposed the desire of innovation on some existing strong groups.Therefore, after the collapse of Qing dynasty, it is natural for them to expect that Yuan Shikai can unite China and realize their dream of innovation and development achieved in a stable process.Based on the needs of authority for people in the trouble time and the traditional culture, the value of order became prominent for intellectuals like Yen fu.On one hand, he knew quite well that to realize the prosperity with constitutional politics at this time, the most important task is to enlighten the public and reforming the society which needs an orderly state under an authoritarian government, and to achieve freedom at this time, it is necessary to have a powerful government which can represents the country to start social governance.The formation and operation of government is the cost people have to pay in order to live an orderly life, besides, the freedom must be obtained in an orderly society (Stan, 2004) [3].He considered that these situations finally made it more difficult to start the social integration and it will be worse for constructing the political center that the constitutional politics need.On the other hand, he has been thinking about the path to keep prosperity under the constitutional politics in his old age.Seeking the prosperity is temporary while the sustainable development and keeping the prosperity is permanent.So it is the core that how to form a rational prosperity-remaining order.After having a meeting for more than two weeks, the advisory council now became powerless.He said that they must resign and they would be dissolved by themselves if their resignations were not allowed (Sun, 2003) [4].In the places of parliamentary democracy where the elites of constitutional politics gathered, the intellectuals neither had the compromising consciousness nor had an authority to command, which led to an incomplete cycling path for political freedom and the gradually depleted resources.Therefore, it is quite difficult to form a prerequisite for the pursuit of prosperity, then it will be unavailable to own an orderly life to remain the prosperity.
Conclusion
Generally speaking, although it is obvious to see the seeking of prosperity value in Yen Fu's early thought of constitutional politics, in fact, the value of freedom and order has already appeared in his mind which are invisible at that time.The value of prosperity is both the strongest voice of the time and the visible aim value of constitutionalism that he in his early life and other intellectuals has proposed.In a series of official calls to arms and enlightening translation books which exhibited his thought of constitutional politics, he creatively explained the value of prosperity, freedom and order.As the first person to input the Western liberal values in modern China, he attached more importance on the prosperous values and began to think about the value of freedom and order in his early time.Then in the mid-term, he returned to the freedom value and talked more about personal freedom and political freedom.Finally, he considered in the late period that order would be the first important task for us at present if the problems of private army regime having appeared after the revolution of 1911 and completing the social reform under the new system and subsequently realizing freedom and eternal prosperity were solved.There are an end wise static value system and a laterally dynamic value relationship in his thought of constitutional politics.In the end wise static value system, the prosperity is the factological value which is properly existing in the situation of saving the nation from extinction; order is the ideological value, which is aiming at the confusion of the early republic of China.The real condition people are living in and the constitutional government; besides, it naturally comes from the needs of the orderly development of the nation and the eternal prosperity.The freedom is the mythological value, which is the ultimate reason why prosperity and freedom are necessary.The value of prosperity, freedom and order in the lateral dynamic value system is dialectical subjects which can mutually promote and are all relying on each other.The unilateral relationship between means and goals only leads to the destruction of three sides.All in all, on one hand, this thesis wishes to be analyzed from his constitutionalism thought to understand the varying features of inner motivation-the change of the aim value with the passing time that owned by the intellectuals who are on the stage of constitutional politics at the end of Qing dynasty from the point to area.In turn, we can also understand the constitutional practice of our predecessors better through the values they have paid attention to.On the other hand, the current practice of constitutional democracy in the process of democratization also requires that society can condense the value consensus to promote the political reform of socialism with Chinese characteristics together.In addition, it needs to find the basic accordance of value ideas from the commendable experimental field of the one hundred constitutional politics which contain untiring efforts of our predecessors. | 3,697.4 | 2015-09-15T00:00:00.000 | [
"Philosophy"
] |
Optimal Kiefer Ordering of Simplex Designs for Third-Degree Mixture Kronecker-Models with three ingredients.
The research was financed by the Deutscher Akademischer Austausch Dienst (DAAD), Abstract This paper investigates the Kiefer optimality in the third-degree Kronecker model for mixture experiments. For mixture models on the simplex, a better design is obtained, by matrix majorization that yields a larger moment matrix due to increase of symmetry and Loewner ordering. The two criteria together constitute the Kiefer design ordering and any such criteria single out one or a few designs that are Kiefer optimal. For the third-degree mixture models with three ingredients, an exchangeable moment matrix was constructed by use of Kronecker product algebra. These moment matrices are symmetrical, balanced, invariant and have homogenous regression entries which are good and have desirable properties for an optimal design. Then, the necessary and sufficient conditions for two exchangeable third-degree K-moment matrices to be comparable in the Loewner matrix ordering were set up. The weights obtained from the original design were used in the construction of the weighted centroid designs. Based on the results obtained, it was shown that the set of the weighted centroid designs constitutes a minimal complete class designs for the Kiefer design ordering and that any design that is not weighted centroid design can be improved upon by convex combination of an appropriate elementary designs. Key words: Kiefer optimality, Kronecker product, weighted centroid designs, simplex centroid design . DOI : 10.7176/MTM/9-1-06
Introduction
The description of the original mixture problem is when two or more ingredients are mixed together to form a product. This product has desirable properties that are of interest to the manufacturers. It is assumed that these properties are functionally related to the product composition and that by varying the composition through the changing of ingredients proportions, the properties of the product will also vary. In the general mixture problem, the measured response is assumed to depend only on the relative proportions of the ingredients present in the mixtures and not on the amount of the mixture Cornell (1990). The study of functional relationship between the measured property (response) and the controllable variables is to determine the best combination of ingredients that yield the desired product. In this basic example of Cake formulations using baking powder, shortening, flour, sugar, eggs and water, the experimenter is looking for fluffiness of the cake, such that fluffiness is related to the ingredient proportions. Similarly, in building construction concrete formed by mixing sand, water, and one or more types of cement building, then the desired property is the hardness or compression strength of the concrete, where the hardness is a function of the percentages of cement, sand, and water in the mix. Cornell (1990) lists numerous examples and provides a thorough discussion of both theory and in practice. Therefore, a mixture experiment involves varying the proportions of two or more ingredients, called components of the mixture, and studying the changes that occur in the measured properties (responses) of the resulting end products. Clearly, if we let q represent the number of ingredients (or constituents) in the system under study and if we represent the proportion of the ith constituent in the mixture by ,then; 0, = 1,2, … , and ∑ = 1+ =1 2+ 2+ … + = 1.0 (2) t ,shown in Eqs.
(1) and (2) the geometric description of the factor space containing the q components consists of all points on or inside the boundaries (vertices, edges, faces, etc.) of a regular (ql)-dimensional simplex. For q = 2 components, the simplex factor space is a straight line. With three components q = 3, the simplex factor space is an equilateral triangle, and for q = 4 the simplex is a tetrahedron.
Mixture Experiments
, that is the factors represent relative proportions of m ingredients blended in a mixture (Cheruiyot et al, 2017). The experimental conditions are points in the probability simplex, which constitute the independent and controlled variables (factors). Replication under identical experimental conditions or responses from distinct experimental conditions are assumed to be of equal (unknown) variance 2 and uncorrelated. The functional relationship between dependent and independent variable within the range of interest is represented by a Taylor polynomial of low degree, d. There are three types of mixture design; simplex-lattice design, simplex-centroid design and simplex axial design. When the mixture components are subject to the constraint that they must sum to one, then standard mixture designs for fitting standard models used are Simplex-Lattice designs and the Simplex-Centroid designs. A simplex design is a mixture design in which the design points are arranged in a uniform way known as lattice. The word, lattice means an array of points and is used in reference to specific Taylor polynomial equation. The simplex centroid design is constructed to form a triangle with data points located at each corner, at the three midpoints on each side and as well as the point located in the centre(centroid). In the simplex lattice design, the points are located on the vertices and mid-edges of an equilateral triangle only, gives more information about response surface behaviour for binary blends, while the points that are located within(inside) the triangle in simplex centroid design and axial design, more uniform distribution is in the interior of the triangle Pukelsheim, 1998a, 1999).
Simplex centroid Designs
The simplex is defined in geometrical terms as a regular figure, where all of its angles are congruent and all of its sides are congruent, such as equilateral (3-sided), tetrahedron (4-triangular faces), and other polygons with triangular faces. In the simplex centroid design, the points are located on the vertices, mid-edges and in the centre (centroid) a triangle .Generally, in a q-component simplex-centroid design, the number of distinct points is 2 − 1.These points correspond to q permutations of (1,0,0, … ,0 ) or q single-component blends, the ( 2 ) permutations of ( 1 2 , 1 2 , 0,0, … ,0 ) or all binary mixtures, the ( 3 )| permutations of ( 1 3 , 1 3 , 1 3 , 0,0, … ,0 ) ,…and so on, with finally the overall centroid point ( 1 , 1 , … , 1 ) or q-nary mixture. In other words, the design consists of every (nonempty) subset of the q components, but only with mixtures in which the components that are present appear in equal proportions. Such mixtures are located at the centroid of the (q-l)-dimensional simplex and at the centroids of all the lower dimensional simplices contained within the (q-l)-dimensional simplex. At the points of the simplex-centroid design, data on the response are collected and a polynomial is fitted that has the same number of terms (or parameters) to be estimated as there are points in the associated design (Muriungi et al, 2017). For example, a case where q = 3 component system and the factor space for all blends is an equilateral triangle, then each component assumes the proportions = 0, 1 2 and 1 for i = 1, 2, 3 (3) Setting d = 2 for the proportions in equation (1), that is the second-degree model is used to represent the response surface over the triangle, then The three points which are defined as (1, 0, 0) or 1 = 1, 2 = 3 = 0; (0,1,0) or 1 = 3 = 0, 2 = 1 and (0,0,1) or 1 = 2 = 0, 3 = 1 (Cornell, 1990(Cornell, , 2002 The points ( (4), then if the second-degree model is used for a three-component system, we have the expected responses and the polynomial equation of the form; = 1 1 + 2 2 + 3 3 + 12 1 2 + 13 1 13 + 23 2 3 + 123 1 2 3 which is a polynomial fitted that has the same number of terms (or parameters to be estimated) as there are points in the associated design where 1, , 2 , 3 , 12 , 13 , 23 , 123 unknown parameters.
Methodology
The Kronecker product has been applied in this study to derive the exchangeable moment matrices since Kiefer design ordering does not depend on the coordinate system that is used to represent the regression function, though both Kronecker and the Scheffe' are based on the same space of regression polynomials, but differ in their choice of representing this space. (Draper and Pukelsheim, 1998) and (Prescott, et. Al, 2002) put forward several advantages of the Kronecker model such as homogeneity of regression terms, attractive symmetry, compact notation, great transparency, and invariance properties. We refer to the corresponding expressions as K-models or K-polynomials. In particular, polynomial regression model for mixture experiments as suggested Pukelsheim, 1998a, 1999) in the first and second-degree Kronecker mixture models in which they obtained the results for Kiefer design ordering of mixture experimental design were reviewed.
For a linear model with regression function f (t), the statistical properties of a design are captured by its moment and its regression function is given by for first, second-and third-degree Kronecker model respectively.
Kronecker product
The Kronecker product, denoted by , is an operation on two matrices of arbitrary size resulting in a block matrix. The Kronecker product should not be confused with the usual matrix multiplication which is an entirely different operation. For a m k matrix A and an n l matrix B, their Kronecker product The Kronecker product approach bases second-degree polynomial regression in m variables such that ) , , , on the matrix of all cross products is given by, rather than reducing them to the Box-Hunter minimal set of monomials ). , , , The benefits are that distinct terms are repeated appropriately, according to the number of times they can arise, so that transformational rules with a conformable matrix R become simple )( ( and that the approach extends to third-degree polynomial regression. However, the arrangement of triple products k j i t t t in a set of "layered" matrices appears rather awkward. This is where Kronecker products prove useful; they achieve the same goal with a more pleasing algebra. The idea underlying the use of Kronecker products is familiar from elementary statistics, that is the Kronecker product of a vector s ℝ m and another vector t ℝ n then simply is a special case, One of the key property of Kronecker product is the product rule . This greatly facilitates our calculations when we now apply Kronecker products to response surface models of third-degree Kmodel. The first-degree K-model which was proposed by Draper, N. R., Pukelsheim, F, (1998)
The first-degree K-model
If the linear model has regression function t t f = ) ( , which is an identity matrix, the statistical properties of a design are captured by its moment matrix, The first-degree moment matrix of an exchangeable design is the with identical on-diagonal entries 2 , the pure second moments and identical off-diagonal entries 11 , the mixed second moments. Furthermore, the simplex restriction entails that is, the entries of any first-degree moment matrix sum to one, for every design on the simplex. For the firstdegree model on the simplex, the regression function is the identity whence the groups and perm (m) coincide, Therefore the support points of must be among the vertices i e . Because of exchangeability the design assigns constant weight m / 1 to each vertex, whence 1 = .Now we view matrix majorization and Loewner ordering together, to obtain the main result on the Kiefer design ordering in first-degree models.
Theorem 1
Among all designs on the simplex , the unique Kiefer optimal design for a first-degree model is the vertex
Proof:
Let be an arbitrary design on the simplex T. Lemma 1 yield . The convex combinations of the vertex points design 1 and of the overall centroid design m exhaust all possible exchangeable first-degree moment matrices.
The second-degree K-model
The second-degree K-model which was proposed by Draper and Pukelsheim (1998b) is of the form An arbitrary design has second-degree K-moment matrix The K-regression function chosen is We call a design with this invariance property an exchangeable design.
(a) Two factors
For the second-degree K-model with two-ingredient, let be an arbitrary exchangeable design on ,then, the second-degree K-moment matrix is of the form; The simplex restriction has the effect on moment matrix, that is the entries of any second-degree K-moment matrix sum to one; and its simplex restrictions entail:
Proof:
For the direct part we assume This means .For the converse part note that for two ingredients, equality of second order moments implies equality of third order moments. The fourth order moment differences then are, using In terms of matrices this means 0 E 1 1 1 1 Again the be vertex points design 1 and the overall centroid design 2 play a special role . Hence the weighted centroid design is well defined. The following theorem joins the partial steps together to obtain the main result on the Kiefer ordering, that the mixtures of the vertex points designs 1 and of the overall centroid design 2 form a minimal complete class.
Theorem 2
In the two-ingredient second-degree model, the set of weighted centroid designs , constitutes a minimal complete class of designs for the Kiefer ordering.
Proof:
Completeness of C means that for every design not in C there is a member in C that is Kiefer better than . That is, we must show that is more informative than , , but that the two are not Kiefer equivalent, .From the above section and with the weights from Lemma 3, the weighted centroid . The implication of the above is that any design which does not consist of a mixture of elementary centroid designs can be improved upon, in terms of symmetry and Loewner ordering, by using an appropriate combination of elementary centroid designs.
(b) Three factors
The second-degree K-moment matrix with three-ingredients is of the form; With an additional moment of mixed order four given by, The simplex restriction has the effect on moment matrix, that is the entries of any second-degree K-moment matrix sum to one; And its simplex restriction entails;
The Third-degree K-model
The third-degree K-model, which was proposed by (Korir, 2008), is of the form; be an arbitrary exchangeable design on , the third-degree K-moment matrix for two-ingredient is of the form; The simplex restriction has the effect on moment matrix, that is the entries of any K-moment matrix sum to one (Korir, et al, 2009) is called a weighted centroid design. In order to find an appropriate set of weights, we equate selected moments of order lower than four: When the lower order moments are expressed using fourth order moments, these weights are seen to be the ones given in the following Lemma.
Lemma 4 Let
be an exchangeable designs on the simplex , with fourth order moments 211 22 31 4 , , , constitutes a minimal complete class of designs for the Kiefer ordering.
Proof:
The Completeness part is as established just as in Theorem 2. For minimal completeness, we remove a weighted centroid from C and assume that . By Lemma 2.7 (Korir,2008), the two designs share the same lower order moments. The latter determine the weights uniquely, contradicting the assumption that and are distinct. Hence the class C is minimal complete. suggested and analysed canonical model forms when the regression function for the expected response is a polynomial of degree one, two, or three. We refer to these as the S-polynomial or S-models. In this paper, the alternative representation of mixture models is used to investigate the third-degree mixture models with three ingredients. This version is based on the Kronecker product algebra of vectors which was introduced by Draper and Pukelsheim (1998Pukelsheim ( , 1999. The Kronecker algebra gives rise to homogeneous model function and moment matrices. We refer to the corresponding expressions as K-models or K-polynomials. In the third-degree mixture model, whereby the S-polynomial and the expected response takes the form;
Design problem
and when the regression function is the homogeneous third-degree K-polynomial, the expected response takes the form (Korir, 2008), (Gregory et al, 2014).
Exchangeability in Third-degree K-model
Given an arbitrary design , we obtain an exchangeable design (permutation invariant) by averaging over the permutation group, Otherwise the average is an improvement over , in that it exhibits more symmetry and balancedness. In terms of matrix majorization, the moment matrix of the average design is majorized by the moment matrix of ,such that ; (Korir,B.C,2008).
This Therefore, a third-degree K-moment matrix is said to be permutationally invariant when then, we speak of an exchangeable thirddegree K-moment matrix. In a third-degree mixture model, the moment matrix is the form; and has all entries homogeneous of degree six and the simplex restriction has an immediate effect on these moment matrices, as follows That is, all the entries of any third-degree K-moment matrix sum to one; for every design on the simplex.
.2 Kiefer Design Ordering
The optimality properties of designs are determined by their moment matrices (Pukelsheim 1993, chapter 5). We compute optimal design for the polynomial fit model, the third degree Kronecker model. This involves searching for the optimum in a set of competing exchangeable moment matrices (Gregory et al, 2014). The Kiefer partial ordering is a two-stage ordering, reflecting an increase in symmetry by matrix majorization and a subsequent enlargement in the Loewner ordering (Pukelsheim, 2006). In view of the initial symmetrization step, it suffices to search for improvement in the Loewner ordering sense, among exchangeable moment matrices only. First, we obtain the exchangeable moment matrices, then find the necessary and sufficient conditions for two exchangeable third-degree K-moment matrices to be comparable in the Loewner matrix ordering. The comparison of moment matrix inequalities reduces to the comparison of individual moment inequalities which is part of the condition. In terms of matrix majorization relative to the congruence action that is induced on the moment matrices by (Kennedy, et al, 2015).
Further, the weights were derived from the initial design and these are assigned to the points of support in experimental domain , these are points on or inside the boundaries (vertices, edges, faces, Centres) of a regular dimensional simplex. These weights were used to obtain the weighted centroid designs in which a convex combination of the elementary centroid designs give rise to the set of weighted centroid designs. Pukelsheim (1993) gives a review of the general design environment. Klein (2002) showed that the class of weighted centroid designs is essentially complete class for m ≥2 for the Kiefer ordering design. As a consequence, the search for optimal designs may be restricted to weighted centroid designs for most criteria.
.Third degree K-model with Three Factors
In the third-degree model, with three-ingredients proposed by Korir, B.C.(2008), an exchangeable moment matrix on design is of the form where A, B, C, D, F, and G are 9 9 block matrices as follows, That is, the entries of any third-degree K-moment matrix sum to one, The set of moments of order six determines all lower order moments. For instance, the pure fifth moments expand to order six by 51 6 3 2 1 In this way we get the following relations: be the vector of moments up to order five as given in lemma 3.1 (Korir,.2008), then the sixth order moment differences are as follows; Lemma 5 Let and be two exchangeable designs on the simplex . Then we have
There are three elementary centroid designs: 1 is supported on the vertices, 2 on the edge midpoints and 3 the overall centroid point The moments of order six of these designs are respectively ( 1 321 1 222 1 411 1 33 1 42 1 ( 3 321 3 222 3 411 3 33 3 42 3 In order to find an appropriate set of weights for In the three-ingredient third-degree model, the set of weighted centroid designs (Korir, 2008) + + = ) , , ( :
Proof:
The Completeness of C means that for every design not in C there is a member in C that is Kiefer better than . That is, we must show that is more informative than , means that .The implication of the above proof is that any design which does not consist of a mixture of elementary centroid designs can be improved upon, in terms of symmetry and Loewner ordering by using an appropriate combination of elementary centroid designs.
Conclusion
This study investigated the Kiefer design ordering in the third-degree Kronecker model for mixture experiments. For mixture models on the simplex, the improvement of a given design is obtained, by increase of symmetry that yields a larger moment matrix under the usual Loewner ordering. The two criteria together constitute the Kiefer design ordering. For the third-degree mixture models, three ingredients, an exchangeable moment matrix was obtained for each factor-case, then the conditions of any two designs to be comparable were set up; by use of moment matrices. The construction of weighted centroid designs becomes visible. The weights were obtained from an original design, which are used in the construction of the weighted centroid designs. It is shown that the set of the weighted centroid designs constitutes a minimal complete class designs for the Kiefer design ordering. It is also shown that any design that is not weighted centroid design can be improved upon by convex combination of an appropriate elementary design. This study agrees with other studies done earlier for the second -degree Kronecker mixture models by Pukelsheim (1998, 1999). The results obtained were used to get the information matrices and therefore Kiefer optimal designs, hence Kiefer optimality (Gregory. K,2012).
Acknowledgement
I am very grateful for the funding availed to me through the Deutscher Akademischer Austausch Dienst (DAAD). This funding facilitated my undertaking of graduate program and this work could not have been completed without such a generous offer. My sincere gratitude goes to the members of the department for the moral support they accorded me. | 5,209.4 | 2019-01-01T00:00:00.000 | [
"Mathematics"
] |
A Machine Learning-Based Diagnostic Model for Crohn’s Disease and Ulcerative Colitis Utilizing Fecal Microbiome Analysis
Recent research has demonstrated the potential of fecal microbiome analysis using machine learning (ML) in the diagnosis of inflammatory bowel disease (IBD), mainly Crohn’s disease (CD) and ulcerative colitis (UC). This study employed the sparse partial least squares discriminant analysis (sPLS-DA) ML technique to develop a robust prediction model for distinguishing among CD, UC, and healthy controls (HCs) based on fecal microbiome data. Using data from multicenter cohorts, we conducted 16S rRNA gene sequencing of fecal samples from patients with CD (n = 671) and UC (n = 114) while forming an HC cohort of 1462 individuals from the Kangbuk Samsung Hospital Healthcare Screening Center. A streamlined pipeline based on HmmUFOTU was used. After a series of filtering steps, 1517 phylotypes and 1846 samples were retained for subsequent analysis. After 100 rounds of downsampling with age, sex, and sample size matching, and division into training and test sets, we constructed two binary prediction models to distinguish between IBD and HC and CD and UC using the training set. The binary prediction models exhibited high accuracy and area under the curve (for differentiating IBD from HC (mean accuracy, 0.950; AUC, 0.992) and CD from UC (mean accuracy, 0.945; AUC, 0.988)), respectively, in the test set. This study underscores the diagnostic potential of an ML model based on sPLS-DA, utilizing fecal microbiome analysis, highlighting its ability to differentiate between IBD and HC and distinguish CD from UC.
Introduction
Ulcerative colitis (UC) and Crohn's disease (CD), which constitute inflammatory bowel disease (IBD), are characterized by chronic inflammation of the intestines [1].The current diagnostic approach for IBD involves a comprehensive strategy that involves medical history, blood and stool analyses, endoscopy with histological findings, and radiological imaging.However, these methods have inherent limitations as they rely on subjective interpretations without any gold standards and must rule out diseases that appear as IBD, leading to inconsistent results [2][3][4][5].Consequently, the complexity of the diagnostic processes and the absence of specific markers often result in a median time to diagnosis of 3.7 months for UC and 8.0 months for CD, with diagnosis delays exceeding 6.7 and 15.2 months for UC and CD, respectively [6][7][8].Unfortunately, the disease can progress rapidly and present acute exacerbation, leading to disease-related complications such as stricturing or penetrating disease, necessitating intestinal surgery [6][7][8].Therefore, timely diagnosis is crucial for the initiation of effective treatment.
Recently, interest in the role of the gut microbiome in IBD pathogenesis has increased [9][10][11].Emerging evidence suggests that alterations in the composition and function of the gut microbiome contribute to the progression and therapeutic response of IBD [12].This potential link between the gut microbiome and IBD underscores the need for innovative diagnostic tools that utilize fecal microbiome analysis as a noninvasive and easily accessible approach.These notions have been reinforced by numerous studies that have identified alterations in microbial diversity and specific bacterial taxa in patients with IBD compared with those in healthy individuals [13][14][15][16].Distinctions between the fecal microbiomes of patients with UC and CD have been reported, suggesting the possibility of a classification based on these differences [13,14].Furthermore, machine learning (ML) models have shown promising performance in distinguishing between patients with IBD and healthy individuals and between UC and CD [14,[17][18][19].These tools may help differentiate between individuals with IBD and those who are healthy and distinguish between UC and CD, two subtypes of IBD.
In contrast to ML algorithms utilized in previous studies, such as random forest (RF), sparse partial least squares discriminant analysis (sPLS-DA) has several advantages.The primary benefit of sPLS-DA is its ability to select a subset of informative variables to discriminate between classes.Additionally, choosing a sparse set of features helps manage many variables that may not contribute meaningfully to the classification task.Moreover, selecting variables with the most discriminative power can contribute to the creation of an interpretable model.
No studies have used sPLS-DA to differentiate between patients with IBD and healthy controls (HCs) or between patients with UC and CD.Therefore, we implemented a prediction model using ML with sPLS-DA to distinguish between both IBD and HC and UC and CD, demonstrating its performance [20].
Research Cohorts and Sample Collection
We enrolled two patient cohorts, one comprising individuals with UC and the other comprising patients with CD, along with a cohort of healthy controls (HCs).The present study was undertaken in parallel with a retrospective multicenter study performed by an IMPACT (identification of the mechanism of CD occurrence and progression through an integrated analysis of both genetic and environmental factors) [21].In 2017, the IMPACT study team was established in Korea and obtained a national grant to organize a retrospective cohort of patients with CD (aged > 8 years) to identify the mechanisms underlying the occurrence and progression of CD.A total of 16 university hospitals are currently participating in this study and collecting clinical data and biological specimens (namely blood, stool, and tissue specimens) from patients with CD who were newly diagnosed or followed up at these institutions.Patients with UC were selected from a prospective multicenter inception cohort study established for UC multi-omics research in Korea in 2020.Fourteen university hospitals participated in this study and collected clinical data and biological specimens, namely blood, stool, tissue, and saliva samples, from patients with UC.Lastly, the HC group consisted of healthy men and women aged 28-78 years who underwent regular health checkups, including body mass index, smoking status, alcohol consumption, and basic blood tests, annually or biennially at the Kangbuk Samsung Healthcare Screening Center from June to September 2014.This cohort comprised individuals who reported the absence of specific diseases using a self-report questionnaire.Further details are provided in a previous study [22].An HC dataset was acquired by communicating with the authors.
Fecal samples were collected by participants (5 g each) and immediately stored in a deep freezer at −80 • C after submission.The collection time for the UC cohort as an inception cohort was the date of research registration before the initiation of drug therapy.Meanwhile, for the CD cohort with a retrospective design, wherein the patients were already diagnosed and were undergoing treatment, the collection times varied.To minimize these effects, fecal samples were collected after more than 3 months of discontinuing antibiotics or probiotics if the patient was taking them.
Sample Preparation and 16S rRNA Gene Sequencing
Information regarding sample preparation and sequencing can be found in a previous report [23].Briefly, the samples were centrifuged at 15,000 rpm for 20 min at 4 • C to separate the cellular pellet from the cell-free supernatant.DNA was extracted from the cellular pellet using a QIAamp DNA Microbiome Kit (Qiagen, Valencia, CA, USA) in accordance with the manufacturer's instructions.
For 16S rRNA amplicon sequencing, we targeted the high-resolution V3-V4 region, which is identical to the existing HC dataset [22] for comparability.The 16S rRNA gene's V3-4 region was amplified with Illumina adapter overhang sequences using 341F (5 ′ -TCG TCG GCA GCG TCA GAT GTG TAT AAG AGA CAG CCT ACG GGN GGC WGC AG-3 ′ ) and 805R (5 ′ -GTC TCG TGG GCT CGG AGA TGT GTA TAA GAG ACA GGA CTA CHV GGG TAT CTA ATC C-3 ′ ) primers.PCR-generated amplicons were purified using a magnetic bead-based system (Agencourt AMPure XP; Beckman Coulter, Brea, CA, USA).Indexed libraries were prepared by limited-cycle PCR using the Nextera technology, cleaned, and pooled at equimolar concentrations.Paired-end sequencing was performed on an Illumina MiSeq platform using a 2 × 300 bp protocol, according to the manufacturer's instructions.
Data Processing and Downstream Analysis
We employed a streamlined pipeline [24] based on HmmUFOtu (version 1.5.1)[25] to analyze the 16S rRNA amplicon sequencing data, as described below.Quality filtering of raw sequence data was performed using fastp [26].Following the recommendations of fastp (version 0.23.2),sequences with a quality score below 20 and reads with a length of less than 150 bp were excluded, as described in a previous study [24] for HC sample processing using fastp.To perform reference-based operational taxonomic unit (OTU) clustering, each trimmed read was individually aligned to the HmmUFOtu model to generate a continuously aligned sequence for each pair.Subsequently, the contig sequences were positioned onto the reference phylogenetic tree (derived from GreenGene version 13.8 and the RDP Classifier Training set version 18) and assigned to the nearest node using the HmmUFOtu main program.The Biostrings (version 2.54.0)Bioconductor package was employed to generate consensus sequences by aggregating the amplicons associated with a shared HmmUFOtu node.We employed Mothur (version 1.48.0)[27] for de novo chimera checking of the consensus sequences, Kraken2 (version 2.1.2) [28] with default parameters, and SILVA reference (version 138.1) for taxonomic assignment.
Microbiome profile data were analyzed using phyloseq (version 1.38.0), a Bioconductor R package.Non-bacterial sequences and those lacking phylum-level annotations were excluded from the analysis.In subsequent analyses, we utilized the cut-off that had yielded significant results in an earlier study [24], excluding samples with fewer than 20,000 read counts and rarely observed phylotypes.We used the Bioconductor R package microbiome (version 1.16.0) to compute the alpha diversity indices for the samples.Using Mothur, we calculated beta diversity indices and conducted permutational multivariate analysis of variance (PERMANOVA) tests based on distance matrices to examine the differences in microbiome composition between different phenotypes.
Machine Learning for Disease Prediction Model
Given the merging of the datasets sequenced at different time points, we used ANCOM-BC (version 1.4.0)[29], specifying the covariate as the time point to adjust for batch effects among the sample groups sequenced at different times before constructing the ML model.We identified the fractions of taxonomic groups with significantly different absolute abundances at each time point.Subsequently, to mitigate variations owing to differences in sequencing depth among samples, we performed a log transformation by adding a pseudo-count of one and subtracting this fraction from the log-transformed abundance obtained from ANCOM-BC.
For subsequent steps, such as principal component analysis (PCA) and prediction model development, we used the mixOmics R package (version 6.18.1) in Bioconductor.We utilized sPLS-DA for variable selection, interpretable results, and computational efficiency.
Because our data were somewhat imbalanced, we matched the age, sex, and number distribution of each class group by downsampling the dataset before training the ML model.The dataset was then randomly divided into 70% training and 30% test sets while maintaining the class proportions.
We employed feature selection and parameter optimization, as recommended by mixOmics.First, we trained the initial sPLS-DA models and assessed their performance with 50 repeated 5-fold cross-validations (5-CVs) to determine the optimal number of components by monitoring the overall error rate trend.Subsequently, we performed tuning processes to select the features for each component.Using these optimal parameters, the final sPLS-DA model was developed, and its performance was measured.
To avoid bias or loss of information, the entire model development process, including downsampling, was repeated 100 times with random shuffling of the training and test splits.Subsequently, the average performance was assessed.
Processing of 16S rRNA Gene Amplicon Sequencing Data
We performed 16S rRNA gene amplicon sequencing of stool samples from 2247 individuals, constituting three phenotypic groups: 671 with CD, 114 with UC, and 1462 HCs.The characteristics of each group are presented in Table 1.
During sequencing, we obtained 164,539,577 paired-end reads.After quality control, 157,961,202 reads remained.Following reference-based OTU clustering, we identified 88,927 OTUs.Taxonomic assignment and phylotyping of the remaining 83,562 OTUs after chimera removal led to the identification of 2525 phylotypes.In the abundance table filtering step, we filtered out phylotypes with abundances less than 10, those that did not belong to bacterial taxa, or those lacking specific phylum information from the entire dataset.Additionally, samples with a total abundance of less than 20,000 were excluded, resulting in 1517 phylotypes and 1846 samples.We used this dataset (CD, n = 670; UC, n = 113; HC, n = 1063) for subsequent analyses.Detailed information regarding each processing step is presented in Table 2.
Diversity Analysis
The results of the alpha diversity analysis showed that the stool microbiome in HC individuals was significantly richer than that in CD (p < 1 × 10 −2 ) and UC (p < 1 × 10− 4 ) patients (Figure 1a,b); however, between CD and UC, the alpha diversity indices were not significantly different.
Beta diversity principal coordinate analysis (PCoA) plots based on Jaccard and thetaYC dissimilarity indices (Figure 1c,d) showed a distinct separation between the IBD and HC samples along the PCoA1 axis, although there were some overlaps.In contrast, the CD and UC samples remained indistinguishable based on components 1 and 2 in both plots.
Multiclass Disease Prediction Model
Before model development, we conducted a log transformation and bias correction of the stool microbiome profile data using ANCOM-BC.To correct for the bias introduced by different sequencing time points in the profile data, we specified the input covariate of ANCOM-BC as the time-point information (seven time points).Taxonomic groups with significantly different absolute abundances at each time point were identified.Subsequently, we added a pseudo-count of one to the profile data, performed a log transformation, and subtracted the fraction obtained from the ANCOM-BC results.
Initially, we employed the sPLS-DA algorithm to create a multiclass ML model.The entire dataset (CD: n = 670, UC: n = 113, HC: n = 1063) was downsampled to match the age, sex distribution, and class counts (CD: n = 113, UC: n = 113, HC: n = 113) and then split into training and test sets with equal class balance.We allocated 70% of the samples to the training set (CD, n = 79; UC, n = 79; HC, n = 79), and 30% were assigned to the test set (CD, n = 34; UC, n = 34; HC, n = 34).This process was repeated 100 times to demonstrate the robustness of the model.In each repetition, the sPLS-DA model of the training set was initialized to identify the optimal components by monitoring the overall error rate.Subsequently, the tuning process selected the best features for each component.We defined
Multiclass Disease Prediction Model
Before model development, we conducted a log transformation and bias correction of the stool microbiome profile data using ANCOM-BC.To correct for the bias introduced by different sequencing time points in the profile data, we specified the input covariate of ANCOM-BC as the time-point information (seven time points).Taxonomic groups with significantly different absolute abundances at each time point were identified.Subsequently, we added a pseudo-count of one to the profile data, performed a log transformation, and subtracted the fraction obtained from the ANCOM-BC results.
Initially, we employed the sPLS-DA algorithm to create a multiclass ML model.The entire dataset (CD: n = 670, UC: n = 113, HC: n = 1063) was downsampled to match the age, sex distribution, and class counts (CD: n = 113, UC: n = 113, HC: n = 113) and then split into training and test sets with equal class balance.We allocated 70% of the samples to the training set (CD, n = 79; UC, n = 79; HC, n = 79), and 30% were assigned to the test set (CD, n = 34; UC, n = 34; HC, n = 34).This process was repeated 100 times to demonstrate the robustness of the model.In each repetition, the sPLS-DA model of the training set was initialized to identify the optimal components by monitoring the overall error rate.Subsequently, the tuning process selected the best features for each component.We defined the final sPLS-DA model for each run using these optimal components and phylotypes, and we evaluated the performance of each model using the corresponding test set.
Overall, these multiclass models showed suboptimal performances in classifying CD and UC, although the HC group was distinctly identified (Table 3 and Figure 2).the final sPLS-DA model for each run using these optimal components and phylotypes, and we evaluated the performance of each model using the corresponding test set.
Overall, these multiclass models showed suboptimal performances in classifying CD and UC, although the HC group was distinctly identified (Table 3 and Figure 2).
Hierarchical Disease Prediction Model
We chose to create two binary prediction models by observing the suboptimal performance of the multiclass model.The first distinguished IBD from HC samples, and the second classified IBD samples as CD or UC.This hierarchical approach enabled accurate classification of the three phenotypes.
Creating a Predictive Model for Distinguishing between IBD and HC
The entire dataset (CD, n = 670; UC, n = 113; HC, n = 1063) was transformed into a binary classification dataset to distinguish between IBD and HC samples.Initially, 113 CD and 226 HC samples were selected and matched for age and sex with the UC samples.Subsequently, the dataset was divided to yield a 70% training set (CD, n = 79; UC, n = 79; HC, n = 158) and a 30% test set (CD, n = 34; UC, n = 34; HC, n = 68).The CD and UC samples in both sets were merged into the IBD class to form training (IBD, n = 158; HC, n = 158) and test sets (IBD, n = 68; HC, n = 68).This process was iterated 100 times using the same ML procedure applied to each split.The model performance was subsequently averaged across splits to provide a comprehensive evaluation.
As shown in Figure 3a, a representative final model produced a plot with a clear distinction between IBD and HC samples.The performance of each model was evaluated
Hierarchical Disease Prediction Model
We chose to create two binary prediction models by observing the suboptimal performance of the multiclass model.The first distinguished IBD from HC samples, and the second classified IBD samples as CD or UC.This hierarchical approach enabled accurate classification of the three phenotypes.
Creating a Predictive Model for Distinguishing between IBD and HC
The entire dataset (CD, n = 670; UC, n = 113; HC, n = 1063) was transformed into a binary classification dataset to distinguish between IBD and HC samples.Initially, 113 CD and 226 HC samples were selected and matched for age and sex with the UC samples.Subsequently, the dataset was divided to yield a 70% training set (CD, n = 79; UC, n = 79; HC, n = 158) and a 30% test set (CD, n = 34; UC, n = 34; HC, n = 68).The CD and UC samples in both sets were merged into the IBD class to form training (IBD, n = 158; HC, n = 158) and test sets (IBD, n = 68; HC, n = 68).This process was iterated 100 times using the same ML procedure applied to each split.The model performance was subsequently averaged across splits to provide a comprehensive evaluation.
As shown in Figure 3a, a representative final model produced a plot with a clear distinction between IBD and HC samples.The performance of each model was evaluated by predicting the disease class of individuals in the corresponding test sets.The IBD versus HC prediction achieved a mean accuracy of 0.950 (0.890-0.993), sensitivity of 0.918 (0.809-0.985), specificity of 0.985 (0.918-1), and precision of 0.984 (0.910-1) (Table 4).
by predicting the disease class of individuals in the corresponding test sets.The IBD versus HC prediction achieved a mean accuracy of 0.950 (0.890-0.993), sensitivity of 0.918 (0.809-0.985), specificity of 0.985 (0.918-1), and precision of 0.984 (0.910-1) (Table 4).We assessed the abundance of the top 10 phylotypes (Table S1) that played key roles in predicting both the IBD and HC groups in the test set using a heatmap.Except for a few samples, we noticed that there was distinct clustering based on class using the 10 phylotype abundance criteria (Figure 3b).
Creating a Predictive Model for Distinguishing between CD and UC
In the original 100 splits mentioned in Section 3.4.1,we exclusively selected CD and UC samples to establish training sets (CD, n = 79; UC, n = 79) and test sets (CD, n = 34; UC, n = 34) to develop models aimed at distinguishing between CD and UC.These datasets were utilized for model development and evaluation, following an earlier procedure.
A representative split sample plot showed a clear separation between the CD and UC samples (Figure 4a), indicating effective differentiation using stool microbiome data.We conducted a performance evaluation of the trained sPLS-DA models by predicting the disease phenotypes of individuals in the test sets.Across the 100 test sets, the classification results displayed a mean accuracy of 0.956, sensitivity of 0.941, specificity of 0.949, precision of 0.950, and AUC of 0.923 (Table 5).These results indicated that the fecal microbiome-based model could distinguish between CD and UC with excellent performance.We assessed the abundance of the top 10 phylotypes (Table S1) that played key roles in predicting both the IBD and HC groups in the test set using a heatmap.Except for a few samples, we noticed that there was distinct clustering based on class using the 10 phylotype abundance criteria (Figure 3b).
Creating a Predictive Model for Distinguishing between CD and UC
In the original 100 splits mentioned in Section 3.4.1,we exclusively selected CD and UC samples to establish training sets (CD, n = 79; UC, n = 79) and test sets (CD, n = 34; UC, n = 34) to develop models aimed at distinguishing between CD and UC.These datasets were utilized for model development and evaluation, following an earlier procedure.
A representative split sample plot showed a clear separation between the CD and UC samples (Figure 4a), indicating effective differentiation using stool microbiome data.We conducted a performance evaluation of the trained sPLS-DA models by predicting the disease phenotypes of individuals in the test sets.Across the 100 test sets, the classification results displayed a mean accuracy of 0.956, sensitivity of 0.941, specificity of 0.949, precision of 0.950, and AUC of 0.923 (Table 5).These results indicated that the fecal microbiome-based model could distinguish between CD and UC with excellent performance.Using a heat map, we examined the abundance of the top 10 phylotypes (Table S2) that were crucial for predicting both the CD and UC groups in the test set.We observed a distinct clustering based on class using the abundance criteria for the ten phylotypes, except for a few samples (Figure 4b).AUC, area under the curve.
Performance Evaluation of Models in Hierarchical Manner
In the previous step, we noted the effectiveness of the fecal microbiome-based binary classification model in distinguishing IBD from HC, CD, and UC.We evaluated the performance of a hierarchical approach that integrates the two models to predict unknown class labels in the input samples.First, the samples were classified as either IBD or HC; then, among those categorized as IBD, further classification into CD or UC was performed.This hierarchical model was evaluated using test sets to assess its effectiveness.Table 6 presents the results obtained from 100 test sets, showing a mean accuracy of 0.936.It also reveals specific values for CD sensitivity of 0.888, CD precision of 0.965, UC sensitivity of 0.933, UC precision of 0.964, HC sensitivity of 0.956, and HC precision of 0.891.Using a heat map, we examined the abundance of the top 10 phylotypes (Table S2) that were crucial for predicting both the CD and UC groups in the test set.We observed a distinct clustering based on class using the abundance criteria for the ten phylotypes, except for a few samples (Figure 4b).
Performance Evaluation of Models in Hierarchical Manner
In the previous step, we noted the effectiveness of the fecal microbiome-based binary classification model in distinguishing IBD from HC, CD, and UC.We evaluated the performance of a hierarchical approach that integrates the two models to predict unknown class labels in the input samples.First, the samples were classified as either IBD or HC; then, among those categorized as IBD, further classification into CD or UC was performed.This hierarchical model was evaluated using test sets to assess its effectiveness.Table 6 presents the results obtained from 100 test sets, showing a mean accuracy of 0.936.It also reveals specific values for CD sensitivity of 0.888, CD precision of 0.965, UC sensitivity of 0.933, UC precision of 0.964, HC sensitivity of 0.956, and HC precision of 0.891.
Discussion
This study demonstrated the effectiveness of an ML model based on sPLS-DA, utilizing fecal microbiome data, in distinguishing between individuals with IBD and HC, as well as in differentiating between CD and UC.First, we constructed a multiclass ML model to differentiate among CD, UC, and HC.It performed well in distinguishing HC from IBD (CD or UC) with a mean sensitivity and precision of 0.952 and 0.814, respectively.However, it performed poorly in differentiating between CD and UC, with a sensitivity and precision <0.5.To overcome this limitation, we restructured two binary classification models in the next step: one to distinguish IBD from HC and the other to distinguish CD from UC.Using binary classification models, the AUC for distinguishing IBD from HC and CD from UC were outstanding, with values of 0.992 and 0.988, respectively.These findings have substantial implications as they demonstrate robust predictive capabilities.
The strength of this study lies in the pioneering use of sPLS-DA to construct a prediction model for distinguishing between IBD and HC, as well as between UC and CD.sPLS-DA method employed in this study offers several advantages over conventional ML approaches for analyzing fecal microbiome data.It effectively addresses challenges related to high-dimensional data and multicollinearity, while providing interpretability [19].We implemented the ML model based on the training sets and initially confirmed its efficacy in distinguishing IBD from HC and UC from CD. Subsequently, we validate its robustness using separate test sets.This study contributes to the growing body of evidence supporting fecal microbiome analysis for diagnosing and distinguishing IBD [13][14][15][16].
Consistent with previous reports, this study found that CD and UC exhibited lower alpha diversity than that of HC [30][31][32].Beta diversity analysis revealed relatively distinct differences in phylotype distribution using the Jaccard dissimilarity metric, although some overlap was observed with the thetaYC dissimilarity metric.The Jaccard dissimilarity metric focuses on the presence or absence of taxa across samples, and it does not consider their abundance or relative abundance.In contrast, the thetaYC dissimilarity metric considers both the presence and absence of taxa and their relative abundances.In summary, both patients with CD and UC exhibited distinct bacterial taxa that differentiated them from HC.Previous research has also reported differences in taxa between UC and CD compared to HC, although the extent of these differences varies [30,33].
In patients with IBD, the predominant characteristics included an increase in the Proteobacteria phylum, Fusobacterium species, and Escherichia coli [31,[33][34][35][36], while there was a decrease in protective taxa such as Faecalibacterium prausnitzii and Bifidobacterium species [32,33,[37][38][39][40].However, information regarding taxonomic differences between IBD and HC varies among the studies conducted thus far, necessitating further clarification regarding the distinctions between UC and CD.Differences among studies, such as sample type, age, sex, dietary habits, disease extent, disease activity, and concomitant therapies, are likely to influence microbial community structure and diversity [30,31,40,41].Therefore, the application of enumerative information as a diagnostic tool may be limited.This study has clinical value in overcoming these limitations and leveraging the advantages of the sPLS-DA algorithm based on differences in phylotypes to construct an ML model and demonstrate its performance.This study also examined the top 10 genera in HC, IBD, CD, and UC.Notably, we observed differences in the major microbiota between CD and UC, which can provide additional information beyond what was previously reported.
Recent advancements in ML models that leverage fecal microbiome data have shown promising results in IBD diagnosis.For example, using OTUs, the RF algorithm achieved notable performance, with an area under the curve (AUC) of 0.80 and an accuracy of 0.72 for distinguishing IBD from non-IBD groups.Additionally, it attained an AUC of 0.92 and an accuracy of 0.83 for distinguishing between UC and CD [14].In another study, various feature selection techniques were employed to construct an RF model, which demonstrated acceptable discrimination in external validation, yielding AUCs of 0.74 and 0.76 for diagnosing UC and CD, respectively [18].Furthermore, a different study developed an RF model using taxonomic profiles at the species level, achieving an AUC of 0.93 and an accuracy of 0.86 for UC diagnosis and an AUC of 0.93 and an accuracy of 0.83 for CD diagnosis [19].However, the limitations of previous studies include the use of a global data platform, which leads to heterogeneity in disease activity, sample collection, and analysis methods [14,18].Matching was not conducted to minimize bias in the selection of the non-IBD group [14,18,19].Additionally, the last study, which employed a multiclass model for various diseases, may have been inappropriate for distinguishing IBD from chronic IBDs [19].Finally, some studies lacked external validation [14,19].In our study, the sPLA-DA model exhibited exceptional performance in diagnosing IBD, with a mean accuracy of 0.950.Additionally, it distinguished between UC and CD with a mean accuracy of 0.945, surpassing the performances of previous studies.These advancements in harnessing fecal microbiome data to develop ML models hold great promise for enhancing the diagnosis of IBD diagnosis.
This study had several limitations.First, confounding factors, such as age, sex, diet, and medication, were not fully controlled.Furthermore, when examining the top genera each for IBD, UC, and CD, no clear commonalities were found compared with the significant microbiota increases reported in previous studies [33].Second, the UC and CD cohorts differed in their characteristics.Patients with UC comprised an inception cohort with fecal samples collected post-diagnosis and pre-treatment, whereas fecal samples of patients with CD were collected at various stages of treatment.To address possible modifiable aspects, we took measures such as discontinuing the use of antibiotics or probiotics before collecting fecal samples.Future studies should consider these factors to enhance our understanding of the microbial diversity in CD and UC.Third, this study lacked external validation, which limited its generalizability.However, during the division of the training and test sets, efforts were made to downsample 100 times with matching age, sex, and sample size.Finally, this study focused on Korean patients diagnosed with IBDs (CD or UC).Microbial communities can vary across geographical regions [42], and the findings of this study may not be directly applicable to other populations.
Conclusions
In summary, this study successfully developed a prediction model using the sPLS-DA algorithm for diagnosing IBD and differentiating between CD and UC compared with HC, demonstrating good performance.We are optimistic that the ML model developed using fecal microbiome data can contribute to the early diagnoses of CD and UC, facilitating prompt and effective treatments guided by its predictions.However, further external validations across different geographical regions are required to confirm the applicability of the developed model.
Supplementary Materials:
The following supporting information can be downloaded at: https: //www.mdpi.com/article/10.3390/microorganisms12010036/s1,Table S1: The top 10 genera chosen with high frequency in the sPLS-DA models for distinguishing between IBD and HC.Table S2: Top 10 genera chosen with high frequency in the sPLS-DA models for distinguishing between CD and UC.Informed Consent Statement: Written informed consent was obtained from the patients for publication of this paper.
Figure 2 .
Figure 2. A PLS projection in the subspace defined by the sPLS−DA model's first two components, developed for multiclass prediction.
Figure 2 .
Figure 2. A PLS projection in the subspace defined by the sPLS−DA model's first two components, developed for multiclass prediction.
Figure 3 .
Figure 3. (a) A PLS projection in the subspace defined by the sPLS−DA model's first two components, developed for discriminating between IBD and HC.(b) A heatmap representing the abundance of high−contributing phylotype features for predicting the IBD and HC groups in the test set.
Figure 3 .
Figure 3. (a) A PLS projection in the subspace defined by the sPLS−DA model's first two components, developed for discriminating between IBD and HC.(b) A heatmap representing the abundance of high−contributing phylotype features for predicting the IBD and HC groups in the test set.
Figure 4 .
Figure 4. (a) A PLS projection in the subspace defined by the sPLS−DA models' first two components, developed for discriminating between CD and UC.(b) A heatmap representing the abundance of high−contributing phylotype features for predicting the CD and UC groups in the test set.
Figure 4 .
Figure 4. (a) A PLS projection in the subspace defined by the sPLS−DA models' first two components, developed for discriminating between CD and UC.(b) A heatmap representing the abundance of high−contributing phylotype features for predicting the CD and UC groups in the test set.
Table 1 .
Baseline demographic and clinical characteristics of participants.
Values are expressed as n (%) unless otherwise specified.SD, standard deviation; BMI, body mass index; GI, gastrointestinal tract.
Table 2 .
Information for each processing step.
Table 3 .
Evaluation metrics from prediction using multiclass models.
Table 3 .
Evaluation metrics from prediction using multiclass models.
Table 4 .
Evaluation metrics from prediction using IBD vs. HC models.
AUC, area under the curve.
Table 4 .
Evaluation metrics from prediction using IBD vs. HC models.
AUC, area under the curve.
Table 5 .
Evaluation metrics from prediction using CD vs. UC models.
Table 6 .
Evaluation metrics calculated in hierarchical manner.
Table 5 .
Evaluation metrics from prediction using CD vs. UC models.
AUC, area under the curve.
Table 6 .
Evaluation metrics calculated in hierarchical manner. | 7,609 | 2023-12-24T00:00:00.000 | [
"Medicine",
"Computer Science"
] |
Structural and Electronic Snapshots during the Transition from a Cu(II) to Cu(I) Metal Center of a Lytic Polysaccharide Monooxygenase by X-ray Photoreduction*
Background: Lytic polysaccharide monooxygenases (LPMOs) exhibit a copper center that binds dioxygen for catalysis. Results: We present LPMO structures from Cu(II) to Cu(I) and analyze the transition with quantum mechanical calculations. Conclusion: Reduction changes the copper coordination state but requires only minor structural and electronic changes. Significance: These structures provide insight into LPMO catalytic activation for further mechanistic studies. Lytic polysaccharide monooxygenases (LPMOs) are a recently discovered class of enzymes that employ a copper-mediated, oxidative mechanism to cleave glycosidic bonds. The LPMO catalytic mechanism likely requires that molecular oxygen first binds to Cu(I), but the oxidation state in many reported LPMO structures is ambiguous, and the changes in the LPMO active site required to accommodate both oxidation states of copper have not been fully elucidated. Here, a diffraction data collection strategy minimizing the deposited x-ray dose was used to solve the crystal structure of a chitin-specific LPMO from Enterococcus faecalis (EfaCBM33A) in the Cu(II)-bound form. Subsequently, the crystalline protein was photoreduced in the x-ray beam, which revealed structural changes associated with the conversion from the initial Cu(II)-oxidized form with two coordinated water molecules, which adopts a trigonal bipyramidal geometry, to a reduced Cu(I) form in a T-shaped geometry with no coordinated water molecules. A comprehensive survey of Cu(II) and Cu(I) structures in the Cambridge Structural Database unambiguously shows that the geometries observed in the least and most reduced structures reflect binding of Cu(II) and Cu(I), respectively. Quantum mechanical calculations of the oxidized and reduced active sites reveal little change in the electronic structure of the active site measured by the active site partial charges. Together with a previous theoretical investigation of a fungal LPMO, this suggests significant functional plasticity in LPMO active sites. Overall, this study provides molecular snapshots along the reduction process to activate the LPMO catalytic machinery and provides a general method for solving LPMO structures in both copper oxidation states.
Lytic polysaccharide monooxygenases (LPMOs) are a recently discovered class of enzymes that employ a copper-mediated, oxidative mechanism to cleave glycosidic bonds. The LPMO catalytic mechanism likely requires that molecular oxygen first binds to Cu(I), but the oxidation state in many reported LPMO structures is ambiguous, and the changes in the LPMO active site required to accommodate both oxidation states of copper have not been fully elucidated. Here, a diffraction data collection strategy minimizing the deposited x-ray dose was used to solve the crystal structure of a chitin-specific LPMO from Enterococcus faecalis (EfaCBM33A) in the Cu(II)-bound form. Subsequently, the crystalline protein was photoreduced in the x-ray beam, which revealed structural changes associated with the conversion from the initial Cu(II)-oxidized form with two coordinated water molecules, which adopts a trigonal bipyramidal geometry, to a reduced Cu(I) form in a T-shaped geometry with no coordinated water molecules. A comprehensive survey of Cu(II) and Cu(I) structures in the Cambridge Structural Database unambiguously shows that the geometries observed in the least and most reduced structures reflect binding of Cu(II) and Cu(I), respectively. Quantum mechanical calculations of the oxidized and reduced active sites reveal little change in the elec-tronic structure of the active site measured by the active site partial charges. Together with a previous theoretical investigation of a fungal LPMO, this suggests significant functional plasticity in LPMO active sites. Overall, this study provides molecular snapshots along the reduction process to activate the LPMO catalytic machinery and provides a general method for solving LPMO structures in both copper oxidation states.
Glycoside hydrolases (GHs) 6 are responsible for significant turnover of recalcitrant polysaccharides such as cellulose, hemicellulose, and chitin in nature and are thus of major importance in the global carbon and nitrogen cycles. GHs are extremely diverse enzymes and have undergone extensive characterization and classification, often driven by their potential utilization in the growing biofuels industry (1)(2)(3)(4)(5). More recently, a new class of enzymes was discovered, classified as lytic polysaccharide monooxygenases (LPMOs), which cleave glycosidic linkages in polysaccharides via a copper-mediated, oxidative mechanism (6 -12). LPMOs represent a new enzyme mechanism for the decomposition of recalcitrant polysaccharides and act synergistically with traditional hydrolytic enzymes (6,(13)(14)(15). Unlike GHs, LPMO action generally does not involve a decrystallization step to detach single polysaccharide chains from their insoluble and often crystalline substrates, a process that requires a substantial amount of thermodynamic work (16 -18). Instead, most LPMOs characterized to date are thought to act directly on surfaces of crystalline polysaccharides (6,19). In this manner, LPMOs are able to synergize with hydrolytic enzymes because they are hypothesized to make chain breaks in crystalline regions that are typically thought to be inaccessible for endoglucanases. Conversely, an LPMO able to cleave soluble substrates was recently discovered, indicating an increasing diversity in substrate specificities of these enzymes (20).
LPMOs were previously classified as family 33 carbohydratebinding modules (CBM33s), which range in origin from bacteria to algae, and family 61 glycoside hydrolases (GH61s), which are of fungal origin. CBM33s mined from genomic data often exhibit modular complexity, whereas GH61s are typically either single module enzymes or are bimodular with a catalytic domain and a family 1 CBM (14), similar to many fungal carbohydrate-active enzymes. Recently, Henrissat and co-workers (5) updated the Carbohydrate-Active Enzyme database and classified CBM33s as Auxiliary Activity 10 (AA10) and GH61s as AA9. Another LPMO family was also recently classified as AA11, which exhibits sequence, structural, and electronic characteristics between that of AA9 and AA10 (12).
The chitin-active LPMO from the Gram-negative chitinolytic bacterium Serratia marcescens, CBP21, was the first LPMO to be biochemically characterized (6,21,22). CBP21 catalysis was shown to be dependent on molecular oxygen, an external electron donor, and the presence of a metal ion cofactor (6), later identified as copper (19). Copper ions have been identified to activate AA10 (19,23), AA9 (7,9,10), and AA11 LPMOs (12). In addition to CBP21, LPMO activity has only been demonstrated for two other CBM33s so far, a celluloseactive CBM33 from Streptomyces coelicolor (CelS2 (8)) and a chitin-active CBM33 from Enterococcus faecalis (EfaCBM33A (24)), the latter of which is the subject of this study. EfaCBM33A is the only LPMO found in the genome of E. faecalis and constitutes, along with a GH family 18 chitinase (EfaChi18A), the chitinolytic machinery of the bacterium. E. faecalis is an opportunistic pathogen, and both EfaCBM33A and EfaChi18A are virulence factors (25,26), suggesting a putative second role for these enzymes beyond biomass depolymerization.
The LPMO active site is constituted by two histidine residues (one of which is the N-terminal residue) that coordinate a copper ion in a motif referred to as the "histidine brace" (6,7,10,27). The copper ion is essential for catalytic activity and is likely to be involved in the activation of molecular oxygen (7,9,24). Soluble products resulting from lytic oxidation have been identified as aldonic acids (6, 8 -10, 28) or as oligomers with an oxidized nonreducing end sugar, i.e. a 4-keto sugar (20,29), indicating differences in enzyme regioselectivity. The oxidation products were recently definitively confirmed by NMR spectroscopy (20), and some progress has recently been made toward understanding regioselectivity (30). To date, aldonic acids are the only products observed for AA10 and AA11 LPMOs, whereas both aldonic acids and 4-keto sugars have been observed for AA9 enzymes.
The catalytic mechanism of AA9 LPMOs, which to date only are found in fungi, was recently examined with density functional theory (DFT) calculations (11). Kim et al. (11) predicted that AA9 LPMOs employ a Cu(II)-oxyl reactive oxygen species for hydrogen abstraction from the substrate, followed by an oxygen-rebound mechanism for substrate hydroxylation. This step will be followed by an elimination reaction, resulting in glycosidic bond cleavage. To activate the LPMO catalytic cycle, the initial dioxygen binding was hypothesized to require reduction of Cu(II) oxidation state of the enzyme to a Cu(I) state, likely mediated by an enzymatic or small molecule reducing agent.
Until recently, there has been a dearth of structural data for metal binding in AA10 LPMOs compared with fungal AA9 LPMOs. Recently, Hemsworth et al. (27) reported the structure of Bacillus amyloliquefaciens CBM33 (BamCBM33), with unknown catalytic activity, binding Cu(I). It was shown that BamCBM33 is stabilized in the presence of copper and that the active site of BamCBM33 with a Cu(I) ion adopts a T-shaped geometry (PDB codes 2YOX and 2YOY). A Cu(II) form of Bam-CBM33 was not crystallized, but x-ray absorption near edge structure (XANES) and EPR spectroscopic methods were used to demonstrate that the enzyme was readily photoreduced during crystallization from the Cu(II) form to a Cu(I) state (27).
In this study, we investigate the active site of EfaCBM33 by progressively photoreducing the catalytic copper from Cu(II) to Cu(I) in the x-ray beam using a data collection minimizing the x-ray dose that is deposited in the sample. During photoreduction, we determine successive structural states by collecting x-ray diffraction data sets on the same crystal. By comparing the structures to known Cu(I) and Cu(II) analogues found in the Cambridge Structural Database (CSD) (31), we ascertain that the obtained structures of the EfaCBM33A unambiguously describe varying oxidation states ranging from Cu(II) to Cu(I). Lastly, we conduct quantum mechanical calculations on an active site model of the Cu(II) and Cu(I) forms of EfaCBM33, which suggest that the electronic structure of the active site remains quite similar as measured by atomic charges. Because the initial reduction of Cu(II) is likely a requirement for LPMO activity, these results offer a structural and electronic picture of how LPMO active sites are preactivated for oxygen binding and subsequent catalysis.
EXPERIMENTAL PROCEDURES
Protein Preparation and Crystallization of EfaCBM33A-EfaCBM33A was expressed and purified as previously described (24). The protein was incubated with 1 mM CuSO 4 for 0.5 h. After soaking, the protein solution was desalted using a 10 DG column (Bio-Rad) and concentrated to 25 mg/ml prior to crystallization. EfaCBM33A crystals were obtained in 20% (w/v) PEG 8000 and 0.1 M HEPES, pH 7.5, by the sitting drop vapor diffusion method as previously described (24). Rodshaped crystals grew to an approximate size of 400 ϫ 53 ϫ 40 m after 2 days of incubation at 20°C. The crystal used for data collection was soaked in the crystallization solution with the addition of 20% PEG 400 as a cryoprotectant for ϳ10 s prior to being plunged into liquid nitrogen.
Diffraction Data Collection and Structure Solution-X-ray diffraction experiments were performed at Beamline ID14-EH1 at the European Synchrotron Radiation Facility, Grenoble, France. Six diffraction data sets were collected using the same single crystal. By utilizing a rod-shaped crystal, monitoring the evolution of UV-visible absorption spectra with x-ray dose (32), and a strategy for helical data collection (33), the radiation dose was minimized, and data of a minimally photoreduced state of EfaCBM33 could be collected (34). A helical data collection consists of defining two points on the crystal along the rotation axis of the goniometer. Although the crystal is rotated over a total 97°angular wedge by 1°steps, it is automatically translated along the rotation axis in between two consecutive rotation steps, thus presenting a fresh part of the crystal to the beam for each diffraction frame. Eventually, the x-ray dose deposited in the sample will approximately be d/w smaller than that deposited with a standard data collection protocol, where d is the horizontal distance between the two points, and w is the horizontal width of the beam. Two points on a limited region (320 ϫ 53 ϫ 40 m) of the crystal were set up as start and end points for data collection with a 50 ϫ 100-m x-ray beam. A different exposure dose per image was used for some data sets, as shown in Table 1. Collecting subsequent data sets by this method allowed for the analysis of the effects of photoreduction on the active site copper with minimal systematic errors, because all data sets were collected from multiple and subsequent exposures of the same crystal volume.
All data sets were indexed and integrated using the program XDS (35) and scaled with the CCP4 program suite version 6.2.0 (36). The structures of EfaCBM33A were solved by molecular replacement using CBP21 (22) (PDB code 2BEM) as the starting model. Model building and maximum-likelihood refinement were performed with iterative cycles of model building in COOT version 0.6.2 (37), by inspection of 2mFo-DFc and mFo-DFc A-weighted maps, and model refinement in Refmac5 version 5.6.0117 (38). The bound copper ion was modeled in at a final stage of the refinement. PyMOL 1.5 (39) was used for analysis of the structures and figure preparations. LSQMAN (40) was used for structural alignments, and Ramachandran statistics were determined using MOLEMAN2 (41). Omit maps were calculated using Phenix 1.8.1 (42). Dose calculations for the exposed crystal regions were performed using RADDOSE version 2 (43,44).
Crystallographic Database Search-A search of the Cambridge Structural Database was performed using ConQuest 1.14. Detailed search parameters are described below in the text.
Quantum Mechanical Calculations-Quantum mechanical calculations based on DFT were performed using Gaussian09 (45) on an active site model (ASM) of EfaCBM33, which includes His 29 , Glu 64 , Ala 112 , His 114 , Trp 176 , Ile 178 , Phe 185 , and the copper ion. The crystallographic water molecules that coordinate the copper ion were also included where appropriate. Smaller and larger ASMs were considered (data not shown), and this model was found to provide the optimal balance between reproducing the crystal structure and being computationally tractable. Additionally, as reported under "Results," the RMSDs of the resulting models of both the Cu(II) and Cu(I) states were below 0.4 Å, a value that is well within the range considered to represent sufficient accuracy in cluster models of enzyme active sites (46). All geometry optimizations were conducted with the local meta-GGA M06-L functional (47,48) and the 6 -31G(d) basis set for all atoms. M06-L was chosen because of improved accuracy for the dispersive and mixed binding complexes. The local density function is computationally efficient for optimization of large structures, employing thousands Values for the highest resolution shell are given in parentheses. c Calculated using a strict boundary Ramachandran definition given by Kleywegt and Jones (41).
Transition from a Cu(II) to Cu(I) Metal Center of a LPMO
of basis functions. All ␣and -carbons were fixed during optimizations. All systems were treated with the conductor-like polarizable continuum model (49,50) using diethyl ether solvation (⑀ ϭ 4 for the protein environment) (51). We computed the harmonic vibrational frequencies for all optimized structures to confirm that they are minima, possessing zero imaginary frequencies. Atomic charges were calculated using natural population analysis from NBO 6.0.
RESULTS
Overall Structure of EfaCBM33A in Complex with Copper-The EfaCBM33A with a bound copper atom was crystallized in space group P2 1 2 1 2 1 with cell dimensions of 43.4 ϫ 48.6 ϫ 68.5 Å, one protein molecule per asymmetric unit, and a V m (Matthews coefficient) (53) of 1.97 Å 3 /Da including all C␣ atoms in the structure. We present six structures of EfaCBM33A along the process of x-ray induced photoreduction, all refined at 1.5 Å and final R and R free values of 15.6 -16.1% and 18.3-19.6%, respectively. The data collection and refinement statistics are summarized in Table 1. In all the structure models, there is clear electron density for all the 169 amino acid residues, ϳ285 water molecules, and 1 copper atom bound to the protein. Negligible pairwise RMSD values of 0.03-0.04 Å over all protein atoms show that the structures are essentially identical. The JULY 4, 2014 • VOLUME 289 • NUMBER 27 primary differences are found in the coordination geometry of the copper ion as a function of the x-ray dose. With the exception of the active site, the structure of the copper-bound EfaCBM33A herein is very similar to the previously published apo form without copper (PDB code 4A02 (24); 0.54 Å RMSD on C␣ atoms).
Transition from a Cu(II) to Cu(I) Metal Center of a LPMO
The overall structure and the active site of EfaCBM33A, as observed in the structure determined from the data set obtained after the lowest radiation dose (PDB code 4ALC) exhibits a trigonal bipyramidal (tbp) structure coordinated by two conserved histidine residues (the histidine brace) and two water molecules (Fig. 1, A-C). In this configuration, the N-terminal histidine (His 29 ) forms a bidentate coordination to the copper ion wherein the backbone N atom occupies one of the three equatorial coordination positions, and the side chain N␦ atom occupies one axial position. The other axial position is occupied by the N⑀ atom in the His 114 residue. The remaining equatorial positions are occupied by two water molecules.
Three additional residues conserved in AA10 LPMOs are shown in Fig. 1C, Glu 64 , Ala 112 , and Phe 185 . Ala 112 is not conserved in AA9 LPMOs and is thought to play a role in the potential mechanistic differences between AA9 and AA10 LPMOs (27,54). Phe 185 is located in a similar position to a conserved tyrosine in AA9 LPMOs, which in the latter case therein imparts an octahedral coordination state around the Cu(II) ion (7,13,28,55,56). In AA9 LPMOs, the Glu 64 residue is replaced by a conserved glutamine residue. The 4ALC data set shows two spherical electron densities (1.93 and 1.91 e/Å 3 , respectively) at 2.21 and 2.19 Å from the copper ion, which were modeled and refined as water molecules and are shown as red spheres in Fig. 1C. For the water molecules bound to copper, there are no other stabilizing interactions with the enzyme. Thus, the positions of the water molecules are primarily dictated by coordination to the copper ion.
The only other known AA10 structure with a copper ion bound reported to date is BamCBM33 from Hemsworth et al. . Copper coordination in EfaCBM33A and TauGH61A. Important residues, atoms, and coordination distances to the copper ion in Å. are indicated where appropriate. a, the copper binding site of EfaCBM33A (PDB code 4ALC) displays a trigonal bipyramidal (tbp) coordination of copper and after x-ray exposure adopts a T-shaped (Tsh) configuration (PDB code 4ALT). b, an octahedral Cu(II) coordination in the GH61A from T. aurantiacus (PDB code 2YET). In most AA10 LPMOs, including EfaCBM33A, the tyrosine residue labeled Tyr175 is replaced by phenylalanine. (27), wherein all copper ions were photoreduced to a Cu(I) oxidation state. The BamCBM33 enzyme active site is illustrated in Fig. 1D. The coordination geometry therein is in a T-shaped (Tsh) geometry with no water molecules coordinated to the copper ion. The corresponding protein-copper interactions retain the structure of the histidine brace. The difference in observed geometry between the 4ALC structure and the Bam-CBM33 structure indicates a difference in copper oxidation state, as described in detail further below.
Structural Changes Induced by X-ray Photoreduction-The structural changes caused by the increase in x-ray dosage during photoreduction were limited to the local environment of the copper ion (Fig. 2). Omit map analysis of the copper-coordinated water molecules shows a continuous decay of electron density correlated with x-ray exposure, and at ϳ1 megagray accumulated radiation (4ALT), both water molecules are completely lost. The electron density for the water molecule closest to Ala 112 is retained slightly longer than the other. The decay of the electron density for the two water molecules coordinated to the copper represents a change in the fraction of Cu(II) to Cu(I) populations between the six structures, as a result of the accumulated radiation dose of the exposed crystal region used for data collection. The structures obtained at higher doses of x-ray radiation reveal a continuous shift in the copper coordination configuration from tbp coordination to Tsh geometry in the structures that lack the copper-bound water molecules (Fig. 2).
LPMO Copper Oxidation State Determination by Analogy to Small Molecule Copper Complexes-To monitor the reduction of the copper ion bound to EfaCBM33A, UV-visible microspectrophotometry was used to record spectral changes of the crystal during x-ray exposure. However, the high background noise and low copper ion concentration in the sample prevented successful application of this method. Thus, to ascertain whether the x-ray-induced changes in the conformation of the copper site are indicative of an actual reduction of the copper ion from Cu(II) to Cu(I), the CSD was searched for relevant copper structures (31). The copper coordination in the retrieved structures was then compared with the initial and final EfaCBM33A structures.
The most obvious structural change upon x-ray exposure is that the two water molecules coordinating to the copper ion gradually disappear, as shown in Fig. 2. This demonstrates that the coordination number for the copper ion drops from five to three; the conformation of the copper site changes from a fivecoordinated tbp structure to a three-coordinated Tsh geometry (Fig. 3). Gradual disappearance of electron density upon increasing the x-ray dosage was not observed for any other water molecule in the structure, suggesting that the effects seen for the copper-bound waters relates to a change in the copper ion. Among more than 40,000 copper structures in the CSD (31), nearly half fit our initial search criteria (1 Յ coordination number Յ 8; only nitrogen and/or oxygen as coordinating atoms), including 9,727 five-coordinate and 564 three-coordinate structures (data not shown). Limiting the search to only include those with histidine-like coordination surroundings and excluding strained structures left 10 five-coordinate structures, all Cu(II) ( Table 2), and 24 three-coordinate structures, all Cu(I) ( Table 3). Details regarding their Cu-N/O bonds, bond angles, torsion angles, and the resulting overall geometry of the copper site are shown in Tables 2 and 3, respectively.
The EfaCBM33 structures have two axial copper-nitrogen (Cu-N ax and Cu-N axЈ ) bond distances, which both decrease by Structures found in the CSD criteria described in the text listed by their code, oxidation state, bond lengths [Cu-N, unless noted], bond angle distribution, and configuration. All structures include a five-coordinated copper(II) ion, including NIDHOY (which is mislabelled in the CSD). One structure shows a copper configuration that is very similar to the configuration in 4ALC: ZUBHOT (tbp configuration with two equatorial oxygen atoms with the longest bond distances to the oxygen from the copper ion), whereas CISHIW also has the same general layout (highlighted in peach). When comparing individual bond distances, it should be noted that both ZUBHOT and CISHIW coordinate to one anionic species each, ClO 4 Ϫ and NCS Ϫ , respectively. sqpy, square pyramidal; tbp, trigonal bipyramidal; dist., distorted. The 4ALC structural details are included for reference at the end of the table (highlighted in turquoise).
Transition from a Cu(II) to Cu(I) Metal Center of a LPMO
0.05 Å going from tbp to Tsh, whereas the equatorial coppernitrogen (Cu-N eq ) bond distance becomes 0.075 Å longer. Additionally, the nearly linear N ax -Cu-N axЈ angle in tbp, 176.2°, bends a bit off-axis in Tsh, 167.5°, whereas the N eq -Cu-N ax and N eq -Cu-N axЈ angles increase by 5. (Table 4). Taken together, these observations show that the structural changes observed upon irradiation of EfaCBM33A reflect photoreduction of Cu(II) to Cu(I).
Quantum Mechanical Calculations of the LPMO Active Site-The structures presented above enable DFT calculations to quantify how the electronic structure of the active site changes upon reduction. The active sites of 4ALC, the Cu(II) structure, and 4ALT, the Cu(I) structure, were both examined with the M06-L functional and the 6 -31G(d) basis set, by employing an ASM representation of the system. Quantum mechanical geometry optimizations were conducted with a range of ASMs. The model consisting of the residues His 29 , Glu 64 , Ala 112 , His 114 , Trp 176 , Ile 178 , and Phe 185 was found to yield the smallest RMSD values for a size that was still computationally tractable with a full quantum mechanical treatment of the ASM in both structures (Table 5; RMSDs of 0.37 and 0.32 Å for 4ALC and 4ALT, respectively). Fig. 4 shows comparisons between the crystal structures and the quantum mechanically optimized ASMs. All computed distances between coordinating nitrogen atoms and the copper differ from the crystallographically observed distances by less than 0.07 Å, which is well within the resolution of the structure (Table 5).
Subsequent to the geometry optimizations, natural population analysis was conducted to examine the charge distributions for both states. As shown in Table 5, the copper ion charges in the oxidized and reduced states of the active site are ϩ1.48 and ϩ0.99, respectively. These values agree well with the formal oxidation states of Cu(II) and Cu(I) and also agree remarkably well with the charges found in both the formal Cu(II) and Cu(I) oxidation states of ϩ1.48 and ϩ0.92, respec-
Copper structures with similar coordination as EfCBM33A after being subjected to radiation
Copper structures with similar coordination as EfCBM33A after being subjected to radiation, using the criterion described in the text as listed in CSD by their code, oxidation state, bond lengths, maximum N-Cu-N bond angle deviation (from 120°), maximum Cu-N-N-N torsion angle, and assigned configuration. All structures include three-coordinate copper(I) ions. One structures has nearly identical configuration to the irradiated form of EfCBM33A (4ALT): PIVNOX (Tsh, uneven bond distribution), whereas GUVLUF also has the same general layout (highlighted in peach). trig, trigonal; tpy, trigonal pyramidal; Tsh, T-shaped. pt1 and pt2 denote part 1 and part 2 of the same reported structure, denoting crystallographically independent atoms. The 4ALT structural details are included for reference at the end of the table (highlighted in turquoise).
tively, in an AA9 LPMO with a similar ASM approach (11). Interestingly, the charge distribution of the coordinating histidine residues does not show a significant change, despite the substantial change in the copper ion oxidation state. This result suggests that the LPMO active site is able to readily accommodate both oxidation states of copper with little overall change in the charge distribution in the enzyme.
DISCUSSION
The present study presents the second structure of a CBM33 with copper bound and the first structure of a CBM33 with a Cu(II) ion. Using a data collection strategy allowing for the structure determination of LPMO structures in both copper oxidation states, we were able to visualize structural and copper coordination changes associated with reduction. This experimental methodology is quite generalizable and can be used to capture the electronic and structural transitions in metalloenzyme reduction at advanced light sources.
It has been proposed that an electron from cellobiose dehydrogenase or from a small molecule reducing agent such as ascorbic acid can reduce the LPMO copper ion to a formal oxidation state of Cu(I) (10), prior to binding of dioxygen. This order of events is in accordance with the notion that molecular oxygen tends to bind copper proteins when the metal ion is in the reduced monovalent state (59). Subsequent to dioxygen binding, the catalytic cycle is initiated, which results in substrate hydroxylation, followed by elimination to cleave the glycosidic linkage (10,11). This general mechanism will incorporate a single oxygen atom from molecular oxygen into the products. The elimination product can then undergo a hydrolysis reaction, which will incorporate an oxygen atom from water, as demonstrated in the mass spectrometry experiments performed by Vaaje-Kolstad et al. (6) with CBP21 and 18 Ocontaining reagents. Interestingly, the series of structures
Five-coordinated, dihydrate copper structures in the CSD
Entries selected that exhibit two coordinating water molecules are listed by their CSD code, oxidation state, bond distances, and overall configuration. Notably, the mean Cu-N bond distance is shorter than the mean Cu-O distance. Two copper sites exhibit sqpy configuration (highlighted in red), but most are in a tbp configuration, as also observed in 4ALC. sqpy, square pyramidal; tbp, trigonal bipyramidal; dist., distorted; ax, axial; eq, equatorial. The residues in gray (carbon), blue (nitrogen), and red (oxygen) represent the crystal structures, and the residues shown in green represent the geometry optimized structures from the DFT calculations. The copper is colored gold and green, respectively. The water molecules from the crystal structure are shown as red spheres in 4ALC, and the optimized water molecules are shown in stick format.
Transition from a Cu(II) to Cu(I) Metal Center of a LPMO
described in the present study show very little structural variation in the conformation of the copper site, despite the change in the copper coordination state induced by x-ray photoreduction. This result highlights that the LPMO catalytic center is preorganized to readily accommodate both oxidation states of copper.
Building on the present results, including the use of the CSD to "annotate" copper site configurations, we conducted a survey of previously reported LPMO structures in terms of their copper oxidation state, the results of which are reported in Table 6. The structures that are annotated as having a Cu(II) metal, NcrPMO-2, NcrPMO-3, and TauGH61A, all exhibit an octahedral six-coordinated octahedral binding motif, which is compatible with the copper being Cu(II). Although the variation in bond lengths and angles is quite high between noncrystallographic symmetry-related molecules in some structures, as well as between the different structures, it is reasonable to conclude that all these published structures represent oxidized Cu(II), in accordance with annotation in the PDB. It is interesting to note that known AA9 LPMO structures with copper contain Cu(II), even though specific precautions to prevent x-ray photoreduction do not seem to have been taken. Under standard x-ray conditions, photoreduction of copper bound to CBM33s readily takes place (27), which could indicate a difference between the AA9 and AA10 LPMO copper sites, possibly caused by the extra coordinating tyrosine in AA9 LPMOs (Fig. 3b). Indeed, based on observed structural and EPR-spectrum differences, Hemsworth et al. (27) have suggested that the oxidative chem-istries catalyzed by these enzyme families may differ. Further work is needed to substantiate this hypothesis. Regardless, considering the large overall similarity of the copper sites (Fig. 3), including the histidine brace, it seems reasonable to hypothesize that both enzyme types employ similar catalytic activation steps for reduction of the copper atom to prime the active site for binding molecular oxygen. Thus, it is likely that the structural and electronic insight obtained here for an AA10 LPMO will be relatively similar for an AA9 LPMO enzyme.
Lastly, the DFT calculations employed here reveal that primarily a coordination number change with only very minor geometry changes in the coordinating atoms is required for binding to the copper ion as it goes from a Cu(II) state to Cu(I). As measured by the atomic charges, very little change occurs in the surrounding protein residues electronically. This result is similar to that found for an AA9 LPMO in our recent mechanistic study (11). Therein, we computed the partial charges of the Thermoauscus aurantiacus AA9 LPMO (TauGH61A) upon copper reduction from Cu(II) to Cu(I) as the step before oxygen binding. This calculation showed that the partial charges of the coordinating atoms only very slightly change upon reduction and concomitant removal of the coordinating water molecules, which in an AA9 LPMO changes the copper coordination from distorted octahedral to tetrahedral coordination (11). Taken together, these results suggest that both fungal and nonfungal LPMO active sites are quite plastic and can readily bind both states of copper. Moreover, the development of a robust ASM for AA10 LPMOs will likely enable the study of the complete reaction mechanism of this family of LPMOs using a cluster model or theozyme approach, similar to that done for AA9 LPMOs (11).
CONCLUSIONS
In this study, we present a crystallographic and computational study of the effects of copper reduction in EfaCBM33, using the structures of well characterized small molecule copper complexes from the CSD to assign the oxidation state the copper ion. X-ray photoreduction causes clear changes in the active site of EfaCBM33, namely the loss of the coordinating water molecules. By correlating the structural data with the CSD, the two forms of EfaCBM33A were assigned as a Cu(II) and Cu(I) state with a trigonal bipyramidal and T-shaped geometry, respectively. DFT calculations reveal only minor changes in the atomic charges required for binding to either oxidation state of the copper ion, similar to what was found in a theoretical study for an AA9 LPMO (11). This study provides the first experimental data set to provide insight in the reductive step that activates an LPMO for catalysis. | 7,471.8 | 2014-05-14T00:00:00.000 | [
"Chemistry"
] |
Quantum random number cloud platform
Randomness lays the foundation for information security. Quantum random number generation based on various quantum principles has been proposed to provide true randomness in the last two decades. We integrate four different types of quantum random number generators on the Alibaba Cloud servers to enhance cybersecurity. Post-processing modules are integrated into the quantum platform to extract true random numbers. We employ improved authentication protocols where original pseudo-random numbers are replaced with quantum ones. Users from the Alibaba Cloud, such as Ant Financial and Smart Access Gateway, request random numbers from the quantum platform for various cryptographic tasks. For cloud services demanding the highest security, such as Alipay at Ant Financial, we combine the random numbers from four quantum devices by XOR the outputs to enhance practical security. The quantum platform has been continuously run for more than a year.
INTRODUCTION
Random numbers play an important role in cybersecurity, cryptography, lottery, and scientific simulations [1][2][3] .In recent years, with the widespread of next-generation information technologies such as big data, cloud computing, and Internet of Things, a large amount of confidential data related to customer privacy has been increasingly exposed to the Internet.Data security is facing great challenges from increasing computing power, future quantum computers, and newalgorithm attacks 4,5 .In the meantime, poor implementations of randomness generation would open up serious security loopholes for cryptosystems even when the underlying algorithms are secure [6][7][8][9][10] .Even for those newly proposed lattice or hashed-based quantum-safe cryptography algorithms, randomness is still a fundamental problem that cannot be solved with classical means.The ability to provide high-quality, high-speed, and stable random number services is an essential demand for information security today 11,12 .
Quantum random number generators (QRNGs) have attracted extensive interests in the past two decades.For a review of the subject, please refer to the recent review articles 13,14 and references therein.The essential difference between QRNGs and classical ones (such as pseudo or thermal noise-based) lies in the unpredictability 15,16 .Guaranteed by the principle of quantum mechanics, QRNGs can avoid predictability loopholes in classical random numbers.As a result, quantum devices show the superiority in tasks with a high information security level, such as data encryption, authentications, and digital signatures.Till now, various methods have been applied to generate quantum randomness, such as detecting the path of a single-photon after a beam splitter [17][18][19] , the arrival time of a weak coherent state [20][21][22][23][24] , the photon-counting detection or the vacuum-fluctuations of an optical field [25][26][27][28][29] , and the phase fluctuations in spontaneous emission [30][31][32][33] .Moreover, when we relax the assumptions and characterizations on the devices, there are also device-independent 34,35 and semi-deviceindependent QRNG schemes 36,37 .
With so many different choices of QRNGs, it is challenging for end-users to understand the underlying principles and to get familiar with various physical and application programming interfaces (APIs) of different devices.Besides, no universal QRNG standards and verification techniques have been officially released so far, making it difficult to evaluate the quality and performance of QRNGs.Individual QRNG devices are usually lack of real-time randomness check, and cannot provide sustainable random number services to online security applications with high stability request.A high-quality quantum random number service should be adaptive with various QRNGs using different interfaces and plug-and-play even if any (not all) device fails.
In this work, we realize a platform on the Alibaba Cloud servers that provides random numbers from four different types of QRNGs, including those based on single-photon detection, photon-counting detection, phase-fluctuations, and vacuum-fluctuations.Real-time post-processing and randomness monitoring modules are integrated into the platform.The generated random numbers are fed into applications either on the Alibaba Cloud servers or remote access for data encryption, with various security levels and speeds.For applications in financial services requiring the highest security, we combine the random numbers from the four quantum devices by bitwise exclusive-OR of the outputs.In this case, as long as at least one of the devices provides true randomness, the applications are secure.A universal trust-cloud-center is more reliable than individual device manufactures.In practice, it is much more challenging for hackers to find loopholes in all different QRNGs.In the future, we would add more quantum entropy sources into the systems to further enhance the security on the implementation level.
Platform realization
Recently, popular QRNG realizations are mainly based on singlephoton detection, photon-counting detection, and phase or vacuum-fluctuations.The schematic diagram of each principle is shown in Fig. 1.The most straightforward idea of QRNG is based on single-photon detection, as shown in Fig. 1(a).When a photon passes through a 50/50 beam splitter, the probabilities to enter detector "0" and detector "1" are balanced [17][18][19] .Due to the dead time of single-photon detectors, such QRNGs usually have a limited speed of Mbps, while phase-fluctuation ones [30][31][32][33] in Fig. 1(b) can dramatically increase the generation speed up to Gbps using traditional photodetectors.Due to the complexity of their optical setups, commercial phase-fluctuation QRNGs are normally bulky (approximately in the size of 1U rack).A more compact QRNG-chip (Fig. 1(c)) based on photon-counting detection 25 has been demonstrated and commercialized, with a relative simple setup and moderate generation rate (240 Mbps).However, the theoretical evaluation of classical and quantum entropy for direct photoncounting QRNGs is still under discussion.For comparison and easy demonstration, a lab-made vacuum-fluctuation QRNG [26][27][28] has also been demonstrated in Fig. 1(d) with a generation rate of 400 Mbps (limited by the characteristics of the homodyne detector), where the lab-made homodyne detector has a bandwidth of 150 MHz and a common mode rejection ratio of 30 dB.All the types of QRNGs above are adopted by our platform.
We present the QRNG Platform protocol and the schematic diagram of the platform setup is shown in Fig. 2. Details of the QRNG cloud platform protocol are described as follows.
1. Data import.The cloud platform adopts quantum random numbers from various QRNG devices through different interfaces (e.g., PCIe, USB, Ethernet, etc.).Our online random number server provides standard interfaces (RESTful or gRPC API), or random numbers can be downloaded from website directly for end-users.The request size of random numbers can be customized, and APIs are compatible with multiple data format including binary, text, ASCII, etc. 2. Randomness extraction.The randomness of the input random numbers from different entropy sources are evaluated.The random numbers pass a real-time randomness extractor, by which the randomness per bit is enhanced to almost 1. 3. Bitwise XOR.A bitwise XOR operation is performed between random numbers from two or more quantum entropy sources.For each train of n-bit random series X i (n), the output random train Y(n) is given by This step is optional.Here are some remarks for the protocol.First, we integrate a real-time post-processing into our cloud.The post-processing technique, namely a randomness extractor (k, ϵ, n, d, m)-extractor, is a function that transforms the n-bit raw sequence with conditional minentropy H min (ρ A |E) ⩾ k (this quantity characterizes the true randomness of ρ A in the presence of eavesdroppers E) to a m-bit sequence arbitrarily close to a uniform distribution with the help of a d-random-bit seed.This process will succeed with a probability no <1 − ϵ.The relations of these parameters are given by the Leftover Hash Lemma For different QRNGs, we quantify the randomness and apply corresponding extractors according to H min (ρ A |E) and get uniform output sequences.In our implementation, both commercial and lab-made QRNG devices are connected to the QRNG Platform.For the commercial QRNGs, the conditional min-entropy is evaluated internally with real-time post-processing by the devices.For the lab-made QRNG, we assume that the quantum signal and classical noise follow independent Gaussian distributions 29 in the strong local oscillator limit, and we have the following relation about their variances, σ t is the standard deviation of the output of the ADC, which includes both σ q (quantum signal) and σ c (classical noise).In a QRNG whose devices are trusted and characterized, the conditional-min entropy is calculated by extracting the quantum signal from classical noise.
Here we take a vacuum-fluctuation QRNG as an example.The fundamental quantum randomness comes from the shot noise of the coherent laser source, whose variance σ 2 q is a linear function of the intensity of the local oscillator 38 .The classical noises are assumed to be independent with laser power [26][27][28][29] , which can be obtained in the absence of the local oscillator.Then we can calculate the signal-to-noise ratio at a certain laser power, and the output randomness is given by the min-entropy function where J is the label of ADC bins, Δ is the resolution of ADC, Gð0; σ t ffiffiffiffiffiffiffiffiffiffiffiffiffiffiffiffiffiffiffi γ=ðγ þ 1Þ p Þ is a Gaussian distribution with zero mean and a variance of σ 2 t γ=ðγ þ 1Þ, and is the probability to generate a certain sequence of random numbers.P max is the maximum probability of some random number sequence occurs per sample, which can be calculated by the area under the probability density in an ADC bin.The conditional min-entropy of other type of QRNGs can be calculated with similar process.
Second, the bitwise XOR operation enhances the reliability of the random numbers in case some of the entropy sources are infiltrated by eavesdroppers.According to Shannon entropy theory 1 , in Eq. ( 1), Y(n) has the perfect secrecy to be all possible n-bit train, providing that any one of X i (n) is random and X i (n) trains are independent from each other.If only two trains X 1 (n) and X 2 (n) are applied, Eq. ( 1) is similar to the one-time pad, where X 1 (n), X 2 (n) and Y(n) are corresponding to the key, the plaintext and the ciphertext respectively.Therefore, we do not need to trust all of the QRNGs but only at least one of them.As long as one of the QRNG entropy sources is reliable, the output Y(n) is random.On the other hand, according to the Leftover Hash Lemma Eq. ( 3), the parameter ϵ characterizes the failure probability of the hashing function.Since the integrated QRNGs work independently, the total failure probability can be decreased by the XOR operation according to the union bound.
Third, depending on specific circumstances, different strategies can be applied to meet the requirements of different levels of security and speeds.For example, financial services such as Alipay require the utmost security.We need to close any possible loopholes in the cryptosystem.For this purpose, random numbers from various quantum devices are taken and processed as in Eq. ( 1).As a result, the highest speed is limited by the slowest QRNG at a rate of 16 Mbps.If end-users have concerns with some specific entropy sources or if any of the hardware breaks down, they can always choose an arbitrary combination of these QRNGs.
Finally, end-users are permitted to store and manage the random number files on the cloud.The generated random numbers could be used for encryption required by other services, either on the end-users' (e.g., Ant Financial) or remote-users' servers (e.g., Smart Access Gateway (SAG)).As those services are on the cloud, high-volume random numbers are required by thousands of servers and real-time post-processing is needed to meet the requirements of online encryption.
Practical implementation and applications
Our platform provides high-quality random numbers in a distributed network environment.The generated random numbers can be further combined with encryption protocols, such as Internet protocol security (IPsec) and SSL/TLS.In these protocols, the existing pseudo-random numbers used in key exchanges, authentication, and digital signatures are replaced with quantum random numbers.Pseudo-random numbers generated by deterministic algorithms will inevitably be predictable and reproducible.The quality of pseudorandom numbers is related to the complexity of the algorithm.With the increasing computing power, the security guaranteed by the complexity of the algorithm is seriously threatened.In contrast, QRNGs with intrinsic unpredictability can be used to greatly enhance the security of cryptosystems.
One example of the QRNG service for practical implementations is the SAG data encryption scenario.SAG is a cloud access solution for connecting hardware and software to the nearest Alibaba Cloud resources through the Internet in encrypted mode.SAG can connect branches (or outlets) and local data centers to the cloud, which enables enterprises to access the cloud more intelligently, safely and reliably.Since more than one million enterprise users from multiple industries are connected by SAG, secure access to the cloud is extremely important.The highly-random QRNG service can be used to enhance the security of the SAG, as shown in Fig. 4. The cloud QRNG transmits the highly-random numbers through RESTful (Representational State Transfer) or gRPC (google Remote Procedure Call) API through TLS to the Cloud Console, OpenSSL uses both Diffie-Hellman and PQC algorithms in the public key infrastructure (PKI) and the self-signed certificate authority (CA).The QRNG APIs is also integrated into an OpenSSL engine as a dynamic library in the TLS transmission, replacing all the random number modules inside the OpenSSL.A concrete example is AliVPN, which is a self-defined protocol for data encryption.Random numbers from the QRNG platform are used in AliVPN to enhance randomness of the keys.All the quantumsafe techniques are optional here in order to be compliant with security standards in some circumstances.
Similar implementations have also been demonstrated in Ant Financial, where the QRNGs are integrated into an OpenSSL engine as a dynamic library in data encryption of the Alipay cloud CA center.Quantum and pseudo-random numbers are switchable in the applications, which are compatible with current security standards.The generated random numbers are also distributed to end-users directly for certificate of authority and encryption in TLS protocols.
The QRNG service is worldwide, provided to ten different places at Shanghai, Japan, Hongkong, Singapore, Malaysia, Indonesia, Australia, United Kingdom, Germany, US East, US West, as shown in Fig. 4.
DISCUSSION
True random numbers are critical components in all cryptosystems.The major advantage of the QRNG platform over the other ones is it can avoid the loopholes of predictable random numbers.The motivation to reduce costs and increase robustness in quantum cryptography remains a great challenge, but the demonstrated feasibility of implementing quantum random numbers in cryptosystems represents an important step toward enhancing the security of classical communications using quantum technologies.The applications in SAG and Ant Financial show the practical implementation of quantum technology in data encryption.Our platform demonstrates quantum random number services with sufficient and adaptive generation speeds, reasonably low costs, controllable risks, high stability, and simple maintenance.
Our scheme shows the feasibility of providing high-quality random numbers in a distributed network environment.The random numbers generated by this scheme can be combined with encryption-related protocols (IPsec, SSL/TLS), identity authentication technologies, or key management systems.The cloud QRNG platform can also be accessed by different end-users in QKD systems, and the generated quantum random numbers can be used as seeds during the QKD communication.For future work, we will consider applications with postquantum algorithms and QKD, since the current distribution of random numbers using classical SSL/TLS is still an issue from the quantum-safe point of view.Integrated QRNG-chip embedded into the SAG devices is also under development to meet certain requirements.Finally, we will develop and integrate more different QRNGs to enhance the security and speed of the system.
Performance of different entropy sources
The cloud-based high-performance QRNG platform is compatible with different types of QRNGs, whose randomness depends on various techniques.Different entropy sources can be chosen to generate the final random keys, which helps prevent from instability or randomness issues caused by individual QRNG device, and increases the reliability of the whole system.Online randomness testings have been performed regularly to ensure the quality of the entropy sources by taking advantages of the computing power on the cloud server.The unpredictability of quantum random numbers comes from the basic principles of quantum mechanics, which guarantees the security of encryption.End-users do not need to understand the underlying hardware equipment and related interfaces, and can simply obtain stable, high-speed, high-quality quantum random numbers for data encryption.
As mentioned earlier, the four different types of QRNGs connected to the platform are based on single-photon detection, photon-counting detection, phase-fluctuations, and vacuum-fluctuations with different interfaces and speeds.To ensure the reliability of the platform, a labmade vacuum-fluctuation QRNG device is implemented together with three commercially available QRNG devices.The type of the singlephoton detection QRNG is Quantis-PCIe-16M from ID Quantique, the type of the phase-fluctuation QRNG is QRG-100E from QuantumCTeck, and the type of the photon-counting QRNG is QRN-16 from Micro-Photon Devices.Table 1 shows the random number generators of different types, speeds and interfaces.
4 .
Randomness test.The platform performs regular (upon request, hourly by default) real-time entropy estimation test (NIST SP 800-90B) to evaluate the non-IID entropy of the quantum random sources, as well as standard NIST randomness test (NIST-800-22) to verify the quality and status of the generated random numbers. 5. Identity authentication.The cloud server performs identity authentication with the end-users upon requests, using preshared key (PSK).6.Data download.End-users download random numbers in plaintext or ciphertext with classical encryption protocols (such as secure socket layer and transport layer security, SSL/TLS) according to the their needs.
Fig. 2
Fig. 2 Schematic Diagram of the QRNG Cloud Platform.Redundant backup server and QRNG devices are provided in different server rooms in case of system corruptions.
Fig. 3
Fig. 3 QRNG in quantum-safe VPNs.Data flow of the generated random numbers from Alibaba QRNG Cloud Platform in the Smart Access Gateway.
Standard NIST randomness tests have been performed on the generated quantum random numbers at a size of 1 Gbit (1000 of 1 Mbit) from different sources and the test results are shown in Table2.Note that hundreds of tests have been performed and Table2only shows one typical example of each devices.The test results (p values and proportions) vary for different sets of random numbers.It turns out that the randomness in most of the generated random number is sufficient to pass the NIST tests.The operation of XOR normally helps to improve values of P-VAL and Proportion, together with the decreasing of the total failure probability ϵ of the hashing function in Eq. (3).
Table 1 .
Quantum random number generation from different entropy sources.
Table 2 .
Random numbers from QRNG sources pass the NIST randomness tests.For tests with multiple outcomes, a Kolmogorov-Smirnov (KS) uniformity test has been performed for p values (P-VAL), and the corresponding proportions (Prop.)are averaged. | 4,209.8 | 2021-07-07T00:00:00.000 | [
"Computer Science",
"Physics"
] |
General results for the Marshall and Olkin ’ s family of distributions
Abstract Marshall and Olkin (1997) introduced an interesting method of adding a parameter to a wellestablished distribution. However, they did not investigate general mathematical properties of their family of distributions. We provide for this family of distributions general expansions for the density function, explicit expressions for the moments and moments of the order statistics. Several especial models are investigated. We discuss estimation of the model parameters. An application to a real data set is presented for illustrative purposes.
INTRODUCTION
Adding parameters to a well-established distri bution is a time honored device for obtaining more flexible new families of distributions.(Marshall and Olkin 2007) introduced an interesting method of adding a new parameter to an existing distribution.The resulting distribution, known as Marshall-Olkin (M-O) extended distribution, includes the baseline distribution as a special case and gives more flexibility to model various types of data.The M-O family of distributions is also known as the proportional odds family (proportional odds model) or family with tilt parameter (Marshall and Olkin 2007).
Let F(x) =1-F(x) denote the survival function of a continuous random variable X which depends on a parameter vector β = (β 1 ,..., β q ) T of dimension q.Then, the corresponding M-O extended distribution has survival function defined by where α > 0 and α = 1-α.For α = 1, G(x) = F(x).Marshall and Olkin (1997) have noted that the method has a stability property, i.e., if the method is applied twice, nothing new is obtained in the second time around.
Additionally, the extended model is geometrically extremely stable.If X i = (i = 1, 2, ...) is a sequence of independent and identically distributed random variables with cdf F(x) and if N has a geometric distribution taken values {1, 2, ...}, then the random variables U = min {X 1 , ..., X N } and V = max{X 1 , ..., X N } are distributed as in (1).It implies that the new distribution is geometrically extremely stable.Marshall and Olkin (2007) have called the additional shape parameter "tilt parameter", since the hazard rate of the new family is shifted below (α ≥ 1) or above (0 < α ≤ 1) the hazard rate of the underlying distribution, that is, for all x ≥ 0, h(x) ≤ r(x) when α ≥ 1, and h(x) ≥ r(x) when 0 < α ≤ 1, where h(x) denotes the hazard rate of the transformed distribution and r(x) is that of the original distribution.Some special cases discussed in the literature include the M-O extensions of the Weibull distribution (Ghitany et al. 2005, Zhang andXie 2007), Pareto distribution (Ghitany 2005), gamma distribution (Ristic´ et al. 2007), Lomax distribution (Ghitany et al. 2007) and linear failure-rate distribution (Ghitany and Kotz 2007).More recently, Gómez-Déniz (2010) presented a new generalization of the geometric distribution using the M-O scheme.Economou and Caroni (2007) showed that the M-O extended distributions have a proportional odds property and Caroni (2010) presented some Monte Carlo simulations considering hypothesis testing on the parameter α for the extended Weibull distribution.Maximum likelihood estimation in M-O family is given in Lam and Leung (2001) and Gupta and Peng (2009).Gupta et al. (2010) compared this family and the original distribution with respect to some stochastic orderings and also investigate thoroughly the monotonicity of the failure rate of the resulting distribution when the baseline distribution is taken as Weibull.Nanda and Das (2012) investigated the tilt parameter of the M-O extended family.
The probability density function (pdf) of the M-O extended-F distribution, say g(x), is given by where f(x) = dF(x) / dx is the baseline density function corresponding to F(x).Here after, we refer to the family (2) as the M-O extended-F distribution.General mathematical properties of the M-O extended-F distribution were not derived by Marshall and Olkin (1997) such as moments and moments of order statistics.In this article, we derive some general structural properties of the M-O extended-F distribution including: (i) expansions for the pdf; (ii) general expressions for the moments; (iii) moments of order statistics; (iv) Rényi entropy.We propose several M-O extended-F distributions taken as baseline in the definitions the Weibull, Fréchet, Pareto, generalized exponential, Kumaraswamy and power function distributions.We discuss maximum likelihood estimation of the model parameters.
The article is organized as follows.Section 2 presents expansions for the density function and for the density function of the order statistics.Explicit expressions for the moments and moments of the order statistics of the M-O extended-F distribution are given in Section 3. Rényi entropy is derived in Section 4. Estimation of the model parameters by maximum likelihood is discussed in Section 5. Section 6 presents an alternative method to estimate the model parameters.In Section 7, we propose several M-O extended-F distributions and discuss some of their properties.Simulation results are performed in Section 8.An application of the current family to a real data set is explored in Section 9. Finally, some concluding remarks are presented in Section 10.
(2) MARSHALL AND OLKIN`S FAMILY OF DISTRIBUTIONS EXPANSIONS Consider the series representation which is valid for |z| < 1 and k > 0, where Г(•) is the gamma function.If α 2 (0,1) using ( 3) in (2), we obtain where The density function (2) can be expressed as Hence, for α > 1, using (3) in the last equation yields We now give the pdf of the ith order statistic X i:n , say g i:n (x), in a random sample of size n from the M-O extended-F distribution.The pdf of X i:n can be expressed as If α 2 (0,1) using expansion (??) in the last equation, we obtain where For α > 1, we write 1 -α F(x) = α{1-(α-1) F(x)/ α} and using (Ref: exp), g i:n (x) becomes where Equations ( 4)-( 7) reveal that the density functions of the M-O extended-F distribution and of its order statistics can be expressed as the baseline density f(x) multiplied by an infinite power series of F(x).They play an important role and will be used to obtain explicit expressions for the moments of the M-O extended-F distribution and of its order statistics in a general framework and for special models.
MOMENTS
Here after, suppose that X has the density function (2).We derive general expressions for the moments of X and its order statistics in terms of the probability weighted moments (PWMs) of the F distribution.The PWMs, first proposed by Greenwood et al. (1979), are expectations of certain functions of a random variable whose mean exists.A general theory for these moments covers the summarization and description of theoretical probability distributions and observed data samples, non parametric estimation of the underlying distribution of an observed sample, estimation of parameters, quantiles of probability distributions and hypothesis tests.The PWMs method can generally be used for estimating parameters of a distribution whose inverse form cannot be expressed explicitly.The PWMs for the baseline F distribution are formally defined by Thus, from equations (4) and (5), the sth moment of X for α 2(0,1) and α >1 can be written as Now, using equations ( 6) and ( 7), we can determine the sth moment of the ith order statistic X i:n in a random sample of size n from X for α 2 (0,1) and α >1 as respectively, where the quantities w j,k , v j , u j,l,k and c j,l are defined in Section 2. Thus, the moments of X and X i:n are obtained in terms of infinite weighted sums of PWMs of the baseline F distribution.
RéNYI ENTROPY
The entropy of a random variable is a measure of uncertainty variation and has been used in various situations in science and engineering.The Rényi entropy is defined by (9) MARSHALL AND OLKIN`S FAMILY OF DISTRIBUTIONS where δ > 0 and δ ≠ 1.For furthers details, the reader is referred to Song (2001).For α 2 (0,1), using expansion (3), we can write For α > 1, we obtain Thus, the Rényi entropy of X can be obtained for α 2 (0,1) and α > 1 as , respectively, where An interesting quantity based on the Rényi entropy is defined by S g = -2d I R (δ) / dδ| δ=1 .It is a location and scale-free positive functional and measures the intrinsic shape of a distribution (see, Song 2001).
MAXIMUM LIKELIHOOD
The model parameters of the M-O extended-F distribution can be estimated by maximum likelihood.Let x = (x 1 , ..., x n ) ┬ be a random sample of size n from X with unknown parameter vector θ = (α, β ┬ ) ┬ , where β = (β 1 ,..., β q ) ┬ corresponds to the parameter vector of the baseline distribution.The log-likelihood function By taking the partial derivatives of the log-likelihood function with respect to α and β, we obtain the components of the score vector Setting these equations to zero, U θ = 0, and solving them simultaneously yields the maximum likelihood estimate ┬ .These equations cannot be solved analytically and statistical software can be used to solve them numerically.For example, the BFGS method (see, Nocedal andWright 1999, Press et al. 2007) with analytical derivatives can be used for maximizing the log-likelihood function ℓ (θ).
The normal approximation for the θ can be used for constructing approximate confidence intervals and confidence regions for the parameters α and β.Under conditions that are fulfilled for the parameters in the interior of the parameter space, we have n √ (θ ^ -θ) α N q+1 (0, K(θ) -1 ), where α means approximately distributed and K(θ) is the unit expected information matrix given by whose elements are given by The asymptotic behavior remains valid if K(θ) = lim n →∞ n -1 J n (θ), where J n (θ) is is the observed information matrix, it is replaced by the average sample information matrix evaluated at θ ^, i.e. n -1 J n (θ ^).The observed information matrix is given by MARSHALL AND OLKIN`S FAMILY OF DISTRIBUTIONS whose elements are We can easily check if the fit of the M-O extended-F model is statistically "superior" to a fit using the F model by testing the null hypothesis H 0 : α≠1.For testing H 0 : α≠1 the likelihood ratio (LR) statistic is given by w = 2{ℓ(α ^, β ^) -ℓ(α ^, β ^)}, where α ^ and β ^ are the unrestricted MLEs obtained from the maximization of ℓ under H 1 and β .The limiting distribution of this statistic is x 1 2 under the null hypothesis.The null hypothesis is rejected if the test statistic exceeds the upper 100(1-γ)% quantile of the x 1 2 distribution.
ESTIMATION-TYPE METHOD OF MOMENTS
We now present an alternative method to estimate the model parameters.Since the moments cannot be obtained in closed form, the estimation by the method of moments is complicated.However, after some algebra, we obtain Thus, we can use (12) to construct a new method of estimation, i.e., if x 1 , ..., x 2 is a random sample with survival function (1), we can estimate the model parameters from the equation In Section, we apply the two methods (maximum likelihood and estimation-type method of moments) to estimate the model parameters of the M-O extended family.
SPECIAL M-O EXTENDED MODELS
We motivate the study of Marshall and Olkin's distributions by considering some special models to illustrate the applicability of the previous results.Here, we obtain the moments and Rényi entropy for some special M-O extended-F distributions when the baseline F distribution follows the Weibull, Fréchet, Pareto, generalized exponential, Kumaraswamy and power distributions.Some others M-O extended-F distributions could be proposed and our general results applied to them.Clearly, the quantities ¿ p,r are determined from the baseline F cdf.
The M-O-EW distribution was studied by Ghitany et al. (2005); see also Barreto-Souza et al. (2011).We obtain where the last equation holds for (δ -1)(° -1) ˃ -1 From these quantities, we immediately obtain explicit expressions for the moments, moments of the order statistics and Rényi entropy.If ° = 1, the results correspond to the M-O extended exponential distribution.
M-O EXTENDED FRéCHET DISTRIBUTION
Here, we consider the Fréchet distribution (for x, σ, λ ˃ 0) with cdf and pdf given by F(x)= e -(σ/x) ° and f(x) = λ σ λ x -(λ +1) e -(σ +1) λ , respectively.The pdf and survival function of the M-O extended Fréchet (M-O-EF) distribution (for x ˃ 0) reduce to respectively.After some algebra, we obtain which is valid for p < λ.Applying this result in ( 8) and ( 9), it follows simple expressions for the moments and moments of the order statistics of the M-O-EF distribution.An expression for the Rényi entropy of this distribution is obtained by inserting in ( 10) and ( 11).
M-O EXTENDED GENERALIZED EXPONENTIAL DISTRIBUTION
The pdf and cdf of the generalized exponential distribution, introduced by Gupta and Kundu (1999), for x ˃ 0 and λ, ° ˃ 0, are given by f δ (x) = ° λe -(λx) (1 -e -λx ) ° -] 1 and F(x) = (1 -e -λx ) °, respectively.By replacing these quantities in (1) and ( 2), the pdf and survival function of the M-O extended generalized exponential (M-O-EGE) distribution reduce respectively to (for x ˃ 0) From (4) and ( 5), we obtain the moment generating function (mgf) of the M-O-EGE distribution as for α 2 (0,1) and , BARRETO-SOUZA, ARTUR J. LEMONTE and GAUSS M. CORDEIRO for α ˃ 1.Both formulas hold for t < λ min {1,1+°}.Hence, their moments can be obtained from the derivatives of the mgf at t = 0. We also have which leads to a simple expression for the Rényi entropy.
M-O EXTENDED POWER DISTRIBUTION
Our final special case concentrates on the M-O extended power (M-O-EPo) distribution defined (for x 2(0, 1/θ) and θ > 0) by taking F(x) = (θx) k in (1).For x 2(0, 1/θ), the pdf and cdf of the M-O-EPo distribution are given respectively by Their moments and moments of the order statistics can be obtained by setting in ( 8) and ( 9).We also have , MARSHALL AND OLKIN`S FAMILY OF DISTRIBUTIONS valid for δ (1-k) -1 .Replacing the last expression in ( 10) and ( 11), it follows simple expressions for the Rényi entropy.
SIMULATION RESULTS
In what follows, we shall present Monte Carlo simulation results.All the Monte Carlo simulation experiments are performed using the Ox matrix programming language (Doornik 2006).Ox is freely distributed for academic purposes and available at http://www.doornik.com.The number of Monte Carlo replications was R = 20,000.We apply the estimation methods before discussed in order to estimate the model parameters of the M-O extended-F distribution.We adopt the M-O-EE distribution with pdf and hazard rate given by respectively, where α ˃ 0 and λ ˃ 0. According to Marshall and Olkin (1997), this distribution may sometimes be a competitor to the Weibull and gamma distributions.The authors derived several properties of the M-O-EE distribution.For example, they showed that h(x) is decreasing in x for 0 < α ≤ 1 and that h(x) is increasing in 1 and e -λx ≤ G(x) ≤ e -λx/α for α ≥ 1. Rao et al. ( 2009) developed a reliability test plan for acceptance/rejection of a lot of products submitted for inspection with lifetimes governed by the M-O-EE distribution.
For a random sample of size n from this distribution, the total log-likelihood function for the parameter vector θ = (α, λ) ┬ is given by The components of the score vector x i e -λxi 1 -αe -λxi .
Setting U θ = 0 and solving them simultaneously yields the MLEs α ^ and λ ^ of α and λ, respectively.The observed information matrix for the parameter vector θ = (α, λ) ┬ is given in the Appendix.The estimates of α ^ and λ ^ , come from the simultaneous solution of the non linear equations The simulation of the M-O-EE independent deviates can be performed using where U ~ U (0,1).The evaluation of point estimation was performed based on the following quantities for each sample size: (i) mean; (ii) relative bias (the relative bias of an estimator (θ ^) of a parameter θ is defined as {E(θ ^θ ^)} / θ, its estimate being obtained by estimating E(θ ^) by Monte Carlo); (iii) root mean squared error, MSE √ , where MSE is the mean squared error estimated from R Monte Carlo replications.The sample sizes are taken as n = 50 and 150.The values of the parameters were set at λ = 0.5 and α = 0.2, 0.6, 1.5, 1.8, 2.0, 2.5, 3.0, 4.0, 5.0, 7.0, 10 and 20.
The point estimates are presented in Tables I and II for n = 50 and n = 150, respectively.From these tables, note that the root mean squared errors of α ^ and α ^ increase with α whereas the root mean squared error of λ ^ decreases with α.The relative bias of λ ^ also decreases as α increases.Additionally, in most of the cases, the relative bias of α ^ is smaller than the relative bias (in absolute value) of α ^ .Also, the relative bias and the root mean squared error of λ ^ is smaller than the relative bias (in absolute value) and the root mean squared error of λ ^ in most of the cases.Also noteworthy is that as the sample size increases, the root mean squared error decreases.In general, the maximum likelihood method yields better estimates of α and λ than the estimation-type method of moments.
Tables III and IV evaluate the overall performance of each of the two different estimators, for each value of n.Each entry in Table III corresponds to a specific estimator and a specific value of n.We obtain what we call the integrated relative bias squared norm (Cribari-Neto and Vasconcellos 2002).This is computed as where the r h 's (h = 1, ..., 12) correspond to the twelve different values of the relative bias of each estimator.Similarly, Table IV III provides a measure of the overall performance of the estimators regarding the bias, and Table IV gives a measure of their overall performance regarding the root mean squared error.In short, the figures in these tables reveal that the maximum likelihood method should be preferred than the estimation-type method of moments in order to estimate the model parameters.
APPLICATION
As Marshall and Olkin (1997) pointed out the M-O-EE distribution considered in the previous section may be an alternative to the Weibull and gamma distributions.In what follows, we shall present an empirical application to a real data set in which the M-O-EE distribution may be preferred than the Weibull and gamma models.We consider the active repair times (hours) for an airborne communication transceiver given in Jørgensen (1982).All the computations were done using the Ox matrix programming language (Doornik 2006). of the estimated density of all fitted models are given in Figure 1.Note that the M-O-EE model provides a better fit than the other models.In Figure 2, we plot the estimated hazard rate for the M-O-EE, Weibull and gamma distributions.Notice that the estimated hazard ratio of the M-O-EE distribution is decreasing and belongs to the interval [λ ^; λ ^/ α ^] = [0.1618;0.3935],expected.Now, we shall apply formal goodness-of-fit tests in order to verify which distribution fits better to these data.We apply the Cramér-von Mises (W*) and Anderson-Darling (A*) statistics.The statistics W* and A* are described in details in Chen and Balakrishnan (1995).In general, the smaller the values of the statistics W* and A*, the better the fit to the data.Let H(x; θ) be the cdf, where the form of H is known but θ (a k-dimensional parameter vector, say) is unknown.
CONCLUDING REMARKS
We study some mathematical properties of the Marshall and Olkin's family of distributions (Marshall and Olkin 1997).This family is defined by adding a parameter to a baseline distribution, giving more flexibility to model various type of data.In the last few years, several authors proposed Marshall-Olkin models to extend well-known distributions (see, for example, Ghitany 2005, Zhang and Xie 2007, Ristic´ et al. 2007, Ghitany et al. 2007, Ghitany and Kotz 2007, Gómez-Déniz 2010).In this article, we derive various structural properties of the Marshall-Olkin extended-F distribution not explored before including simple expansions for the density function and explicit expressions for the moments, moments of the order statistics and Rénvy entropy.Our formulas related with this class of distributions are manageable, and with the use of modern computer resources with analytic and numerical capabilities, may turn into adequate tools comprising the arsenal of applied statisticians.Several Marshall-Olkin extended-F distributions are proposed and some of their mathematical properties are given.We discuss maximum likelihood estimation of the model parameters and propose an alternative estimation method.Monte Carlo simulation experiments are also considered.Finally, an empirical application to a real data set is presented.
Figure 1 -
Figure 1 -Estimated pdf of the Marshall-Olkin extended exponential (M-O-EE), Weibull and gamma models for a real data set.
provides the average root mean squared error expressed as √of each estimator.For each value of n, Table Table V lists the MLEs of the parameters (standard errors in parentheses) and the values of the loglikelihood functions.The M-O-EE distribution yields the highest value of the log-likelihood function.Plots MARSHALL AND OLKIN`S FAMILY OF DISTRIBUTIONS
TABLE I Point estimates for n=50 and different values of α.
MSE: mean squared error; Rel.Bias: relative bias.
To obtain the statistics W* and A*, one can proceed as BARRETO-SOUZA, ARTUR J. LEMONTE and GAUSS M. CORDEIRO
TABLE III Integrated relative bias squared norm.
MLE: maximum likelihood estimate.MME: modified moment estimate.
TABLE II Point estimates for n=50 and different values of α.
MSE: mean squared error; Rel.Bias: relative bias.MARSHALL AND OLKIN`S FAMILY OF DISTRIBUTIONS
TABLE IV Average root mean squared error.
MLEs: maximum likelihood estimate.M-O-EE: Marshall-Orkin extended exponential. | 5,073.6 | 2013-03-01T00:00:00.000 | [
"Mathematics"
] |
Characterization of the p90 Ribosomal S6 Kinase 2 Carboxyl-terminal Domain as a Protein Kinase*
The carboxyl-terminal domain (CTD) of the p90 ribosomal S6 kinases (RSKs) is an important regulatory domain in RSK and a model for kinase regulation of FXXFXF(Y) motifs in AGC kinases. Its properties had not been studied. We reconstituted activation of the CTD inEscherichia coli by co-expression with active ERK2 mitogen-activated protein kinase (MAPK). GST-RSK2-(aa373–740) was phosphorylated in the P-loop (Thr577) by MAPK, accompanied by increased phosphorylation on the hydrophobic motif site, Ser386. Activated GST-RSK2-(aa373–740) phosphorylates synthetic peptides based on Ser386. The peptide RRQLFRGFSFVAK, which was termed CTDtide, was phosphorylated with K m and V max values of ∼140 μm and ∼1 μmol/min/mg, respectively. Residues Leu at p −5 and Arg at p −3 are important for substrate recognition, but a hydrophobic residue at p +4 is not. RSK2 CTD is a much more selective peptide kinase than MAPK-activated protein kinase 2. CTDtide was used to probe regulation of hemagglutinin-tagged RSK proteins immunopurified from epidermal growth factor-stimulated BHK-21 cells. K100A but not K451A RSK2 phosphorylates CTDtide, indicating a requirement for the CTD. RSK2-(aa1–389) phosphorylates the S6 peptide, and this activity is inactivated by S386A mutation, but RSK2-(aa1–389) does not phosphorylate CTDtide. In contrast, RSK2-(aa373–740) containing only the CTD phosphorylates CTDtide robustly. Thus, CTDtide is phosphorylated by the CTD but not the NH2-terminal domain (NTD). Epidermal growth factor activates the CTD and NTD in parallel. Activity of the CTD for peptide phosphorylation correlates with Thr577 phosphorylation. CTDtide activity is constrained in full-length RSK2. Interestingly, mutation of the conserved lysine in the ATP-binding site of the NTD completely eliminates S6 kinase activity, but a similar mutation of the CTD does not completely ablate kinase activity for intramolecular phosphorylation of Ser386, even though it greatly reduces CTDtide activity. The standard lysine mutation used routinely to study kinase functions in vivo may be unsatisfactory when the substrate is intramolecular or in a tight complex.
The RSKs 1 (reviewed in Ref. 1) and the closely related nuclear enzymes MSK1 (2) and MSK2 (also known as Rsk-b (3)) are unusual among protein serine/threonine kinases because they possess two kinase domains. This feature was predicted from analysis of the first RSK cDNA after its isolation from Xenopus laevis (4). There are four human RSK genes encoding RSK1, RSK2, RSK3, and RSK4, 2 and all have the same conserved domain structure (1). The NH 2 -terminal catalytic domain (NTD) of the RSKs is most similar to p70 ribosomal S6 kinases followed by members of the protein kinase B and protein kinase C families. All of these enzymes, including the RSK NTDs, require PDK1 phosphorylation of a conserved serine (Ser 227 in RSK2) in the activation loop for activity (5,6). The NTD once activated is functional for phosphorylation of several physiologic substrates (1), perhaps best demonstrated by the ability of a constitutively active form of RSK, containing only the NTD, to induce metaphase arrest in cleaving Xenopus embryos (7).
The RSK CTD evaluated together with the adjacent COOH terminus is most similar in a standard BLAST search to CaMactivated protein kinases I and II. However, features in the activation loop and the COOH terminus reveal its functional relationship to the single domain MAPKAP kinases 3 such as the MNKs (8) and MAPKAP kinase 2 (9). The RSK and MSK CTDs and the single domain MAPKAP kinases have a threonine residue followed by proline in the kinase activation loop, nine and eight residues from the APE motif. None of the calmodulin-activated protein kinases have threonine followed by proline at this position, a necessary feature for MAPK phosphorylation. Similarly, all of the MAPKAP kinases have a MAPK docking motif in the COOH terminus (10) that we believe to be similar in structural disposition to the calmodulinbinding domain encompassing the ␣R2 helix in CaM kinase I (11). Similar to other MAPKAP kinases, the CTD is activated by MAPK via phosphorylation of a conserved T-P site in its activation loop (12), facilitated by the MAPK docking motif (10). The CTD phosphorylates one known substrate, a conserved Ser (Ser 386 in RSK2) (12) in the linker domain that joins the NTD and the CTD. The majority of this linker corresponds to the carboxyl terminus of p70 ribosomal S6 kinase by alignment and thus belongs to the NTD by inference. Within it are several regulatory sites of phosphorylation that correspond to phosphorylation sites in p70 ribosomal S6 kinase, protein kinase B, and protein kinase C enzymes.
Activation of the RSK NTD is dependent on phosphorylation of the regulatory sites in the linker domain (Ser 386 and putative MAPK sites Thr 365 and Ser 369 ). Ser 369 and Ser 386 may be the most important of these because S369A-and S386A-type mutants of full-length RSK are virtually inactive and unresponsive to agonist stimulation (6,12), whereas T365A-type mutants are nearly wild type (12). Ser 386 lies within a docking motif for PDK1, and Ser 386 phosphorylation is required for PDK1 binding and subsequent activation of the NTD (6). The role of Ser 369 phosphorylation is unknown. Ser 369 is not significantly phosphorylated in truncated RSK2 (amino acids 1-389), which has significant constitutive activity, suggesting that Ser 369 phosphorylation plays a role in the regulation of fulllength RSK. Some conformational states of inactive full-length RSK may sterically inhibit NTD activity independent of Ser 227 phosphorylation in the P-loop because a portion of RSK1 is phosphorylated basally at this site, yet is inactive (12,13). Phosphorylation of Ser 369 , and possibly Ser 386 as well, may contribute to relief of this inhibition in full-length RSK. The Ser 369 kinase(s) are U0126-inhibitable (14), pointing to ERKs1-2 or ERK5 as the upstream kinases for Ser 369 phosphorylation in vivo. Also consistent with this conclusion, phosphorylation of the equivalent serine in avian RSK1 is blocked by deletion of the MAPK docking motif (14). Thus, ERKs are the likely physiologic Ser 369 kinases.
Presently there are multiple models for RSK activation, possibly because of the existence of a multiplicity of activation mechanisms for this key signaling protein. One current model for RSK activation is vectorial (6). In this model, MAPK activates the CTD, which in turn phosphorylates Ser 386 in cis, creating a binding site for PDK1, which in turn phosphorylates Ser 227 , activating the NTD. However, some evidence suggests that NTD activation may not always proceed vectorially from CTD activation. A portion of Ser 222 in RSK1 is phosphorylated basally (12) as already mentioned, obviating the requirement for PDK1 phosphorylation. Furthermore, myristoylated avian RSK1 targeted to the plasma membrane is activated in serumstarved cells independent of evident ERK activation (15). RSK mutants rendered kinase-defective in the CTD (by mutation of the essential lysine) are still activated by growth factors but not as robustly as wild type (16,17).
In comparison with the NTD, much less is known about the CTD. In the current view, the NTD is assumed to be the only domain of the two capable of substrate phosphorylation in trans. The properties of the activated CTD as a protein kinase have never been studied. Our results demonstrate that the isolated, MAPK-activated RSK2 CTD is functional as a protein kinase toward peptide substrates. Furthermore, the CTD but not the NTD portion of the full-length protein selectively phosphorylates the best of these peptides (RRQLFRGFSFVAK), which is referred to herein as CTDtide. This peptide substrate allowed us to probe CTD regulation independently in full-length RSK2.
EXPERIMENTAL PROCEDURES
Materials-The plasmids pET-MEK1 R4F/His 6 ERK2 (18), pMT2-RSK2-(aa1-389) and its S386A mutant (6), and pGEX-5X-MK2-EE (19) were generously provided by Melanie Cobb (University of Texas South-western Medical Center, Dallas, TX), Steen Gammeltoft (Glostrup Hospital, Glostrup, Denmark), and Matthias Gaestel (Max Delbrü ck Centrum Molecular Medicine, Berlin, Germany), respectively. We obtained the murine RSK2 cDNA as pMT2 HA-RSK2 (20) from Christian Bjørbaek (Beth Israel Medical Center, Boston, MA) and have deposited its coding sequence as determined for pKH3-RSK2 (10) as GenBank TM accession number AY083469. The monoclonal antibody to RSK2 Thr(P) 577 (21) was kindly given to us by Paolo Sassone-Corsi (CNRS, Strasbourg, France). The synthetic peptides related to the NH 2 terminus of glycogen synthase were a generous gift of Sir Philip Cohen (University of Dundee, Dundee, UK), and the alcohol dehydrogenase repressor protein 1 (ADR1-g), synapsin, and glycogen synthase peptides were generous gifts from Anthony Means (Duke University, Durham, NC). Goat polyclonal anti-RSK2 (C-19) and the horseradish peroxidaselinked anti-goat antibodies were purchased from Santa Cruz Biotechnology; GammaBind Sepharose plus, glutathione-Sepharose 4B, horseradish peroxidase-linked anti-mouse, and anti-rabbit antibodies were from Amersham Biosciences; epidermal growth factor was from Collaborative Biomedical Products; microcystin LR and PD98059 were from Calbiochem; BHK-21 cells were from American Type Culture Collection; Syntide-2 was from the American Peptide Company; S6 and all RSK peptides were from the University of Virginia Biomolecular Research Facility; and the HA peptide was from the Howard Hughes Medical Institute Peptide Synthesis Facility (Duke University). All of the other reagents and products were from standard sources.
Construction of pAC-pET RSK2 CTD-A PCR-based strategy was used to engineer muRSK2-(aa373-740) (wild type and a K451A kinasedefective mutant) in frame into the BamHI and XhoI sites of pET41b (Novagen). The GST-RSK2 CTD and the kinase-defective mutant were transferred to pACYC184, a plasmid that has a p15A origin of replication, using a strategy suggested to us by Dr. Peter Sheffield (Center for Cell Signaling, University of Virginia Health Sciences Center, Charlottesville, VA). PCR was used to add BclI sites to a ϳ2.3-kb fragment amplified with Pfu polymerase (Promega) from pET41a-RSK2-(aa373-740). This fragment spans the T7 promoter for RNA polymerase, the GST tag, the pET41b multicloning site, RSK2-(aa373-740) with its stop codon, and the T7 terminator. The amplified fragment was cloned into the single BamHI site of pACYC184 (GenBank TM accession number X06403). The construct was verified by sequencing. The plasmids pAC-pET RSK2-(aa373-740) and pAC-pET RSK2-(aa373-740)(K451A) have a chloramphenicol resistance marker. The expressed proteins (exclusive of GST) have a 72-amino acid leader polypeptide (derived from codons within pET41b) that contains His 6 and S tags (Novagen) in addition to thrombin and enterokinase cleavage sites.
Phosphospecific Antibody to RSK2 Ser(P) 386 -The immunogen was [Cys]-Gly-Arg-Phe-Ser(P)-Phe-Val-Ala conjugated to keyhole limpet hemocyanin via the cysteine, and the antisera were produced in rabbits by Research Genetics (Huntsville, AL). The IgG fraction was purified from the production bleed by affinity chromatography on immobilized protein G.
Protein Production and Purification-The vectors encoding GST-and His 6 -tagged proteins were transformed into Escherichia coli BL21 cells with or without the bicistronic vector that encodes constitutively activated MEK1 and wild type ERK2 (18). These cells were grown at 37°C to an A 600 of ϳ0.3 and then induced with 0.5 mM isopropyl--D-thiogalactopyranoside for 6 -8 h at room temperature. The proteins were purified from the cell lysates using either glutathione Sepharose or Ni 2ϩ -nitrilotriacetic acid-agarose (Qiagen) essentially as directed by the manufacturer but with the addition of 2 M lithium chloride washes for the GST-tagged RSK proteins. The purified proteins were quantified using the Bio-Rad protein assay with bovine serum albumin as the standard. Active EE-MAPKAPK2 was made according to Engel et al. (19) using pGEX-5X-MK2-EE.
Kinase Activity and K m Determination-P81 paper assays were used to monitor kinase activity of the purified GST or His 6 -tagged proteins. The kinase assays were performed in 40-l reactions containing (final concentrations) 25 mM Hepes, pH 7.4, 2 mM dithiothreitol, 0.25 mg/ml bovine serum albumin, 10 mM MgCl 2 , a peptide substrate (as indicated in the figure legends), 50 M [␥-32 P]ATP (ϳ4000 cpm/pmol) at 30°C for 13 min or the indicated times in the figure legends. Phosphate incorporation into peptide substrate was determined using P81 phosphocellulose paper as described previously (10).
Cell Culture and Transfection-BHK-21 cells were grown in a humidified incubator at 37°C with 10% CO 2 in Dulbecco's modified Eagle's medium supplemented with 10% newborn calf serum. During transfections, the cells were also supplemented with 100 units/ml of penicillin and 100 g/ml of streptomycin. The cells were transfected with 20 g of DNA using a calcium phosphate profection system (Promega) on 150-mm dishes as described previously (22). Post-transfection (ϳ45 h) the cells were serum-starved for 3 h with or without 50 M PD 98059 and then treated with epidermal growth factor (100 ng/ml) for the times indicated. The cell lysates were made as described previously (22).
Immunoprecipitations and Kinase Activity of HA-tagged Proteins-HA-tagged proteins were immunoprecipitated from cleared lysates with 25 g of 12CA5 antibody and subsequently eluted as described previously (22) except that GammaBind-Sepharose was used, and the secondary antibody was omitted. Kinase activity of 3-l portions of eluted immunoprecipitations was determined by the P81 method in 40-l reactions containing (final concentrations) 25 mM Hepes, pH 7.4, 5 mM -glycerophosphate, pH 7.4, 1.5 mM dithiothreitol, 6 M cAMP-dependent protein kinase inhibitor peptide, 15 mM MgCl 2 , 100 M [␥-32 P]ATP (ϳ4000 cpm/pmol), and either S6 (RRRLSSLRA) or CTDtide (RRQL-FRGFSFVAK) at 100 M. These reactions were initiated with ATP and incubated at 30°C for 13 min. Incorporation was corrected by subtracting the peptide phosphorylation observed from control immunoprecipitations performed in parallel from cells transfected with empty vector and treated equivalently. Specific activity of immunoprecipitation kinase assays (to HA Western signal) was determined from quantitative densitometry using ImageQuant software (Molecular Dynamics) essentially as described (22).
Data Analysis-Representative data are shown from experiments that were repeated at least twice (see legends).
RESULTS AND DISCUSSION
Design of a Plasmid for Reconstitution of Active RSK2 CTD in Bacteria-To address whether the RSK2 CTD is functional for phosphorylation of an exogenous substrate, we felt it necessary to test an active, recombinant protein from E. coli for several reasons. Kinase activity detectable from a CTD recovered from Sf9 or mammalian cells could be ascribed to a contaminant in the preparation. E. coli do not express protein serine/threonine kinases, and protein kinase cascades can be reconstituted in bacteria. In particular, a bicistronic plasmid for inducible expression of untagged active MEK1 together with His 6 -tagged ERK2 was created to produce active ERK2 for crystallization (18). The active MEK1 phosphorylates ERK2 in bacteria, producing doubly phosphorylated, fully active ERK2. Active ERK2 activates RSK2 (reviewed in Ref. 1). We hypothesized that this strategy could be used to produce active recombinant MAPKAP kinases, including the RSK2 CTD, provided the MAPKAP kinase would fold correctly in bacteria. That the RSK2 CTD might fold correctly was suggested by previous work (17), wherein GST-RSK1-(aa386 -752) was expressed in E. coli as a soluble protein and was active as revealed by its ability to autophosphorylate via an intramolecular (concentration-independent) mechanism.
For the above reasons, we engineered muRSK2-(aa373-740) to be expressed as a GST fusion protein (see "Experimental Procedures"). The plasmid pAC-pET RSK2 CTD ( Fig. 1) derived from pACYC184 has a p15A origin of replication and is compatible with the ColE1 replicator in pET-MEK1 R4F/His6 ERK2.
GST-RSK2 CTD Is Activated by ERK2 in E. coli-We used phosphospecific antibodies in Western blots to compare the levels of phosphorylation of Ser 386 and Thr 577 in wild type GST-RSK2 CTD and K451A RSK2 CTD obtained from two conditions: expression alone or expression together with active ERK2 (Fig. 2A). RSK2-(aa373-740) encompasses the Ser 386 site, which is known to be an intramolecular substrate in RSK2 for the CTD. The blots were first probed with anti-RSK as a loading control ( Fig. 2A, top panel) and then stripped and reprobed with anti-Ser(P) 386 ( Fig. 2A, middle panel). No signal was detected for the K451A mutant expressed by itself. Ser 386 was phosphorylated in wild type protein expressed in the absence of active ERK2. Expression with active ERK2 greatly increased Ser 386 phosphorylation for both the mutant and wild type proteins. The specificity of the phosphospecific Ser 386 antibody was verified using RSK2 CTD treated with and without calf intestine alkaline phosphatase (Fig. 2B). Because Ser 386 is phosphorylated by the CTD (12, 23) and is not a MEK or ERK2 substrate, RSK2 CTD enzymatic activity is up-regulated by ERK2 in bacteria.
The blot was stripped and reprobed a third time with a monoclonal antibody (21) to Thr(P) 577 , and the regulatory site was phosphorylated by MAPK in the activation loop of the RSK2 CTD ( Fig. 2A, bottom panel). Neither GST-RSK2 CTD nor the kinase-defective K451A mutant expressed alone was immunoreactive with this antibody. Co-expression with active ERK2 induced a large and easily detectable signal from both active and kinase-defective RSK2 CTD, establishing that active ERK2 produced from the bicistronic plasmid phosphorylates RSK2 CTD on Thr 577 in bacteria.
It was somewhat surprising that the K451A protein was detectably phosphorylated on Ser 386 because this serine is not a substrate for ERK2 or MEK1 R4F. Although it is remotely possible that some portion of the Ser 386 phosphorylation occurring in situ is due to the activating kinases, it is more likely that the K451A mutant retains kinase activity. In kinases that contain the conserved lysine in subdomain II (not all do), mutation of that residue decreases but does not completely eliminate kinase activity. The residual retained activity of these mutated kinases is variable and dependent upon the kinase and the amino acid substitution. For example, ERK2 K52R retains ϳ5% of the activity of wild type ERK2 (24) and still weakly autophosphorylates (25). Our results indicate that the K451A mutant of RSK2 CTD retains reduced activity that is also up-regulated by ERK2. This reaction should be facilitated because it can occur intramolecularly (17), if, as seems almost certain, the unidentified site in that study was Ser 386 .
The residual Ser 386 kinase activity of K451A RSK helps to rationalize the ability of full-length K451A-type RSK mutants to be partially activated by ERK (Refs. 17 and 22 and this study). The Ser 386 site must be phosphorylated or else mutated with phosphomimetic residues to bind and possibly to activate PDK1 (6). This is true for either full-length or truncated RSK. Although RSK2-(aa1-389) is phosphorylated by unidentified Ser 386 kinase(s), the phosphorylation of Ser 386 in full-length RSK is most likely catalyzed by the CTD. (Data showing that EGF can induce phosphorylation of Ser 386 in K451A RSK2 are presented later in the proper context (see Fig. 8D).) Thus, NTD activity in K451A-type RSK mutants can be rationalized in part by the residual intramolecular phosphorylation of these lysine mutants.
Bound ERK2 Is Removed by Lithium Chloride Washes- 1. Plasmid map of pAC-pET RSK2 CTD. A segment from pET41B-muRSK2 CTD from the T7 promoter through to and including the T7 terminator was moved into pACYC184 at the BamHI site. This construct encodes a GST leader polypeptide-RSK2-(aa373-740) of predicted mass 72.4 kDa. The leader polypeptide (derived from pET41B (Novagen)) contains GST, His 6 , and S tags as well as thrombin and enterokinase cleavage sites.
MAPK binds to the docking motif in the carboxyl terminus of RSK, and docking facilitates RSK activation in mammalian cells (10,26). ERK2 co-purified from lysates with the GSTtagged RSK when the two proteins were expressed together in cells (Fig. 3A). The docking motif of RSK contains paired arginines that make ionic bonds to glutamic residues in the common docking domain of ERK2 (27). We tested several options for removing ERK2 from GST-RSK2 CTD while the latter was still bound to glutathione beads. Of these washes, 2 M LiCl released the majority of bound ERK2 from both wild type (Fig. 3A, top panel) and K451A RSK2 CTD (data not shown). Final preparations of both proteins contained only a small, substoichiometric amount of ERK2 that was detectable by Western blotting (Fig. 3A, middle panel) but was not evident in Coomassie-stained gels.
After purification, K451A GST-RSK CTD reproducibly retained more ERK2 than GST-RSK2 CTD (Fig. 3B, bottom panel). The relevant difference between these proteins is the amount of Ser 386 phosphorylation (Fig. 3B, middle panel). This serine is contained within an FXF motif (in RSKs, FSF) that is an ERK docking site in proteins such as ELK1, Lin-1, KSR, and phosphodiesterase (28 -30). This motif mediates interactions with ERK2 that are independent of the D domain MAPK docking site. Binding of ERK2 to FXF occurs at physiologically relevant affinities that are decreased up to 10-fold by disruption of the motif (30). Our data suggest that ERK2 is interacting with the Ser 386 site in GST-RSK2 CTD and that phosphorylation of Ser 386 decreases the affinity of the interaction. The carboxyl-terminal docking domain is the predominant mechanism of ERK-RSK interaction. However, it is plausible that the FSF site in RSKs is contributing to interaction with ERK2. Supporting this, ERK2 binds avidly to phenyl-Superose (31). In other proteins, this motif alone can direct phosphorylation of neighboring ERK sites (28).
RSK2 CTD Phosphorylates Ser 386 Peptide-Having estab-lished that WT-RSK2 CTD was able to phosphorylate Ser 386 in bacteria ( Fig. 2A), we tested RSK2 CTD for enzymatic activity toward an exogenous Ser 386 synthetic peptide (RRQLFRGFS-FVAI) (Fig. 4). WT-RSK2 CTD activated by ERK2 phosphorylated the peptide (solid squares), but WT-RSK2 CTD did not. The residual ERK2 in the preparation should not phosphorylate the peptide because the serine is not followed by a proline. This was confirmed by finding (Fig. 4, crosses) that ERK2, with an activity toward MBP that is 20 times greater than the activity of the wild type RSK toward the Ser 386 peptide, did not phosphorylate the Ser 386 peptide. The residual ERK2 in the CTD preparations can be detected using MBP as the phosphoacceptor (data not shown), but Ser 386 is not an ERK2 substrate. To our knowledge, this is the first proof that the isolated CTD can phosphorylate any substrate in trans.
trans-Phosphorylation Correlates with Thr 577 Phosphorylation-Wild type RSK2 CTD that was not ERK2 activated, and hence not Thr 577 phosphorylated, is able to phosphorylate Ser 386 in bacteria (see above) but does not phosphorylate the Ser 386 synthetic peptide (Fig. 4, open squares). This suggests that the wild type protein is active for intramolecular phosphorylation of Ser 386 but needs Thr 577 phosphorylation and potentially a conformational change to become active toward an exogenous substrate. This conclusion is strongly supported by the finding that K451A RSK2 CTD activated by ERK2 in bacteria phosphorylates the peptide (Fig. 4, inset). Our results show that phosphorylation of RSK2 CTD at the Thr 577 site enhanced intramolecular phosphorylation at Ser 386 as expected and conferred the ability to phosphorylate Ser 386 peptide, which is novel.
A prior indication that the RSK2 CTD might function in trans is contained within a report (23) that first identified phosphorylation of the Ser 386 site in Xenopus RSK. However, that part has been ignored because their experiments had alternative explanations, chiefly coming from use of full-length RSK containing the NTD.
Residues at p Ϫ5 and p Ϫ3 Are Critical for Ser 386 Peptide Phosphorylation-Specificity of protein kinases for physiologic substrates is dictated by several factors. Of these, positioning of the kinase with the substrate and primary sequence preference for the substrate are often paramount. The ability of a kinase to select substrates on the basis of primary sequence is due to pockets on the kinase surface that accommodate specific residues in the substrate, and the selection provided can be more or less stringent.
We tested the RSK2 CTD for stringency in substrate selection by comparing the phosphorylation of several mutant peptides to the parent Ser 386 peptide, each at 0.1 mM (Fig. 5, bottom panel). The sequence similarity of RSK CTD to CaM kinase I and MAPKAP kinase 2 dictated inclusion of the peptide sequence surrounding Ser 386 from p Ϫ5 to p ϩ4. The Ser 386 motif (Fig. 5, top panel) is an exact match to the current consensus for CaM kinase I (BXRXX(S/T)XXXB) (32) and for MAPKAP kinase 2 (XXBXRXXSXX) (33), where B is a subset of hydrophic residues in each case. For MAPKAPK2, the optimal p Ϫ5 hydrophobic residue is bulky (Phe Ͼ Leu Ͼ Val Ͼ Ͼ Ala) and is leucine in the physiologic MAPKAPK2 substrate Hsp27 (LNRQLSS). For CaM kinase I, bulky hydrophobic residues are required at both the p Ϫ5 and p ϩ4 positions (32,34). The Ser 386 site of RSK includes bulky hydrophobic residues at p Ϫ5 and p ϩ4, which could be important. In addition, hydrophobic residues are important at other positions in the consensus sequence of other members of the CaM kinase superfamily to which the CTD is related (35,36). CaM kinase II and phosphorylase kinase both select Phe or another bulky hydrophobic residue strongly at p ϩ1 (36). The "hydrophobic motif " FXXFXF(Y) (37) for PDK1 docking contains phenylalanines that could be important for Ser 386 phosphorylation.
The L381K peptide was not a substrate, suggesting that the p Ϫ5 position is important for recognition by the CTD. Arg 383 at p Ϫ3 was required because the R383K and R383G peptides were compromised or not phosphorylated at all, respectively. These data show that p Ϫ5 and p Ϫ3 are critical for peptide phosphorylation by the CTD and suggest that like MAPKAPK2 the CTD prefers a hydrophobic residue and Arg at these positions. We made leucine substitutions for the phenylalanines. All of these mutants (F382L, F385L, and F387L) were phosphorylated by the RSK2 CTD nearly as well as the parent peptide. Some members of the CaM kinase superfamily (AMPactivated protein kinase (34) and CaM kinase I (32)) strongly prefer a hydrophobic residue at p ϩ4. The RSK CTD resembles CaM kinase II (36) in tolerating Lys replacement at p ϩ4. In the case of full-length RSK, intramolecular phosphorylation of Ser 386 a priori could be due to the combined effects of structural presentation of the site in the linker and to preference for a primary sequence. This will require detailed structure-function studies of the linker. However, the fact that residues in the linker (p Ϫ5 and p Ϫ3) are both critical for peptide phosphorylation and are conserved across RSKs suggests that recognition of Ser 386 for intramolecular phosphorylation will be dependent on the primary sequence in the linker.
As a final caveat, we do not know whether the peptide sequence is an optimal sequence. ERKs autophosphorylate intramolecularly on regulatory tyrosine and threonine sites (38,39) that do not conform to the (S/T)P consensus that is absolutely required for phosphorylation of exogenous substrates. We also do not imply existence of additional physiologic substrates for RSK2 CTD from these data, although our findings make the possibility more plausible.
RSK2 CTD Is a Selective Kinase-We tested ϳ20 peptides that have been used to characterize CaM kinase-related enzymes (Table I and data not shown). For comparison, the peptides in Table I were assayed with an equal amount of the active EE mutant of MAPKAPK2 (19) also produced in E. coli. The RSK linker peptides were MAPKAPK2 substrates, as expected. Peptides 1-4 are related to the NH 2 terminus of glycogen synthase and were used to characterize MAPKAPK2 specificity (33). Peptides 1 and 2 are excellent MAPKAPK2 substrates and were phosphorylated by RSK CTD but were poor substrates in comparison with the RSK linker peptides. Peptides that are excellent substrates for other CaM kinase-related enzymes were also tested and were not appreciably phosphorylated. Syntide 2 is a standard CaM kinase II substrate (35); ADR1-g and synapsin (4 -13) are benchmark substrates for CaM kinase I (40). The AMARA peptide is a substrate for AMP-activated protein kinase (34). None of these peptides were phosphorylated by RSK CTD. These results, with Fig. 5, show that the peptide specificity of RSK2 CTD is similar to MAPKAPK2, but it is distinct and much more selective.
Characterization of the I390K Peptide as CTDtide-The I390K peptide was the best substrate for purified, active RSK2 CTD and is referred to as CTDtide (Fig. 5 and Table I). The apparent K m for this peptide is ϳ140 M. This K m is in the upper range for most protein serine/threonine kinases. The specific activity toward the peptide was 1 mol/min/mg (Fig. 6), indicative of an efficient and properly folded enzyme. For reference, the catalytic subunit of protein kinase A, produced in E. coli, has a K m of ϳ40 M for Leu-Arg-Arg-Ala-Ser-Leu-Gly (Kemptide) and a specific activity of ϳ20 mol/min/mg for phosphate transfer (41). For CTDtide to be useful as a probe for CTD activity, the CTD should phosphorylate the peptide in the full-length protein. The NTD should not phosphorylate the peptide or at least should do so poorly in comparison with CTD.
Specificity of the CTD was tested using constructs created to inactivate the two kinase domains in full-length RSK2 (Fig. 7, A and B). K100A RSK2 is kinase-defective for the NTD; K451A RSK2 is kinase-defective for the CTD. The HA-tagged proteins were immunopurified from EGF-treated cells for assay. The K100A mutant is completely kinase-dead for phosphorylation of S6 peptide but still phosphorylates CTDtide. The K451A mutant is kinase-dead for phosphorylation of CTDtide but not for phosphorylation of S6 peptide (25% of wild type). Note the difference in scales. S6 peptide is a much better substrate for the NTD in full-length RSK than CTDtide is for the CTD. To eliminate the possibility that the NTD could account for a significant proportion of CTDtide activity, we compared intact and a kinase-defective RSK2-(aa1-389) (Fig. 7C). A S386A mutant is as kinase-dead as a S227A mutant of the P-loop because it prevents PDK1 from phosphorylating the P-loop (6). Mutation of S386A caused nearly complete loss of S6 kinase activity but did not significantly affect the small amount of apparent CTDtide activity co-purified in this experiment. This residual incorporation is most likely due to contaminating kinases in the immunopurified RSK. Corrections to CTDtide activity were made from mock immunoprecipitates from empty vector controls, usually around 1-5% of CTDtide activity of wild type or K100A RSK2.
We suspect that the CTDtide activity is constrained in the full-length protein. This would be the expected result if the CTD has evolved to phosphorylate only the Ser 386 site in the linker. Evidence for constraint of CTDtide activity in full-length RSK is FIG. 6. Kinetic analysis of CTDtide phosphorylation by GST-RSK2 CTD. The peptide was used as substrate (25-143 mM) in P81 kinase assays as described under "Experimental Procedures." Active GST-RSK2 CTD was used at 5 mg/ml, and the reactions were run for 5 min. The apparent K m value for the peptide is 140 M, and the specific activity is ϳ1 mol/min/mg (n ϭ 3 Ϯ S.D.).
a Activity was measured using the p81 paper assay (see "Experimental Procedures"). The concentration of enzyme and substrate were 6.25 g/ml and 100 M, respectively. The reactions were run in duplicate for 13 min; 100% incorporation was 6.7 or 7 pmol of phosphate for RSK2 CTD and EE-MAPKAPK 2, respectively.
shown in Fig. 7D, wherein kinase activities of the isolated NTD and CTD kinases are compared with the full-length wild type kinases. The data presented in Fig. 7 are normalized for total RSK protein from the Western signal for HA because we observed differences in level of expression of the separated domains. The CTD alone (RSK2-(aa373-740)) is consistently underexpressed, and the NTD alone (RSK2-(aa1-389)) is overexpressed compared with the full-length proteins (data not shown). Normalization was done as carefully as possible, making preliminary runs to find dilutions that were similar and would be in the linear range of the film after nonextended exposures. The separated CTD domain consistently had higher specific activity when compared with the full-length proteins. In additional experiments, portions of the assayed RSK proteins were also analyzed for reactivity to anti-Thr(P) 577 , and specific activities relative to Thr 577 phosphorylation were calculated (Table II). The enhanced specific activity of the isolated CTD is partially explained by increased stoichiometry of Thr 577 phosphorylation. However, the specific activity of CTD alone, normalizing to Thr 577 phosphorylation, is still 6-fold higher than that of either full-length protein (wild type or K100A). This suggests differences in conformation between the isolated CTD and the CTD in the full-length protein that alter the specific activity. Structures to determine the accessibility of the active site in the isolated CTD versus the CTD in full-length RSK would be of interest.
The I390K peptide can be considered CTDtide because the CTD, but not the NTD, portion of the RSK protein phosphorylates it. CTDtide can be used to assay CTD activity in relation to NTD activity in RSK provided that RSK is purified to remove other kinases, such as MAPKAPK2, that would also phosphorylate it and that proper controls are performed. The specific elution of bound RSK with HA peptide in our experiments may also have helped to reduce the amount of contaminating kinases that may stick to beads but are not elutable with HA peptide.
CTD Activation Is Rapid and Parallels NTD Activation-Previously it has not been possible to determine whether the NTD and CTD kinase activities are congruent (i.e. on together/ off together) or dissociated (i.e. CTD-activated and inactivated, toward S6 (hatched bars), in arbitrary units normalized for RSK protein. In C and D: left ordinate, S6; right ordinate, CTDtide. The data are the averages of duplicates (ranges indicated by error bars) but are representative of six experiments. The CTD but not the NTD is required for phosphorylation of CTDtide. C, S6 versus CTDtide, constitutively active wild type RSK2-(aa1-389), and S386A mutant. S386A mutation reduces S6 but not CTDtide kinase activity; the NTD is selective for S6. (Note difference in scale for apparent CTDtide activity here versus in D.) D, S6 versus CTDtide activity, wild type, and K451A HA-RSK2-(aa373-740). The isolated CTD has 6-fold higher specific activity than CTD in full-length RSK2. HA-tagged RSK2 proteins (indicated) were immunopurified from BHK-21 cells stimulated with epidermal growth factor for 5 min (see "Experimental Procedures"). The procedure included specific elution from the immunoprecipitates with HA peptide. Eluted proteins were assayed with S6 peptide (hatched bars) and CTDtide (RRQLFRGFSFVAK) (solid bars). Specific activity a Activity of HA-RSK proteins, eluted from immunoprecipitates of cells treated for 15 min with EGF (3 l), was measured in duplicate using the p81 paper assay with CTDtide as substrate (see "Experimental Procedures"). The eluted proteins were then analyzed for protein content using western blot analysis (anti-HA), and the results are presented as specific CTDtide activity per the HA signal in arbitrary units.
b Eluted immunoprecipitated proteins from above were then analyzed for the amount of Thr 577 phosphorylation (anti-Thr(P) 577 , and the results are presented as ratio of Thr(P) 577 to HA signal, expressed as a percentage.
c The activity results from the p81 assays above presented as specific CTDtide activity per Thr 577 phosphorylation in arbitrary units.
whereas NTD remains active). We assayed the NTD and CTD activities of HA-RSK2 purified from cells treated with EGF for different times (Fig. 8). We found that activation of the CTD closely paralleled NTD activation for this agonist (Fig. 8A). Maximal activation occurred between 5 and 10 min. Activation was transitory for both domains. EGF also caused a time-dependent activation of CTD activity of K100A RSK2, which lacks S6 kinase activity (Fig. 8B). K100A HA-RSK2 was completely inactive toward S6 peptide, consistent with the demonstrated efficacy of K100A HA-RSK2 as a dominant negative (42). Finally, EGF failed at all times to increase phosphorylation of CTDtide by the full length K451A mutant (Fig. 8C). As noted above, growth factors cause a reduced activation of S6 kinase activity in K451A-type mutants. The K451A mutant is detectably phosphorylated at Ser 386 in unstimulated cells and EGF causes a weak but detectable increase in this phosphorylation (Fig. 8D, compare 0 min to 5 min). This is consistent with data discussed above that were obtained in bacteria. The kinetics of Ser 386 phosphorylation in the wild type, included as a control, show EGF activation at the earliest time point, 2.5 min.
These kinetics are compatible with the vectorial model for RSK activation because intramolecular phosphorylation of Ser 386 would be extremely rapid. Furthermore, PDK1 binding and phosphorylation of Ser 227 is potentially too fast to produce detectably slower NTD activation relative to the CTD. The PDK1 steps may even be preempted. There is evidence (13) for a pool of inactive RSK that is already phosphorylated on the NTD activation loop (12). Activation of this pool of RSK would only require phosphorylation of Ser 386 and/or linker MAPK sites and a conformational change. Because this pool is inactive in unstimulated cells despite phosphorylation of the NTD activation loop, RSK activation must also include relief of intrasteric inhibition (43).
The Future for Other MAPKAP Kinases-Reconstitution of activation of other MAPKAP kinases (MNKs, etc.) along similar lines should be feasible. Although it is possible to generate active mutants that circumvent phosphorylation requirements by truncation or acidic replacements, structural studies of enzyme regulation are most informative with authentic phosphorylated enzyme. This will require generation of additional bicistronic vectors to produce the specific activated MAPKs. Reconstitution of protein cascades in bacteria may also prove to be surprisingly specific. MAPKAP kinase 2 was discovered biochemically as an enzyme that was activated in vitro by ERK (44) but was later shown to be activated by p38␣, MAPK specifically in mammalian cells. Consistent with this, we found that co-expression of MAPKAP kinase 2 with active ERK2 in bacteria caused only modest activation (data not shown), which we found surprising. EE-MAPKAP kinase 2 produced in bacteria, in contrast, was much more active (19) (data not shown). | 8,687 | 2002-08-02T00:00:00.000 | [
"Biology",
"Chemistry"
] |
A Numerical Simulation of a Stationary Solar Field Augmented by Plane Reflectors : Optimum Design Parameters
In this study, a theoretical analysis of a solar field augmented by a fixed reflector placed in the front between the top of the preceding row and the bottom of the succeeding row is presented. An analytical model has been developed and used to estimate the solar irradiation. The analytical model is based on the anisotropic sky model, assuming an infinite length of collector and reflector rows. A simulation has been carried out in order to figure out the behavior of the solar field and to find the optimum design parameters of the solar field leading to a maximum solar energy augmentation. The results obtained are depicted synoptically as a relationship between the solar field design parameters and the latitude angle, and this presentation enables us to determine the optimum design parameters in order to achieve the intended percentage improvement of solar radiation incident on the solar field rows at any location on the Northern hemisphere, which presents the novelty of this research. Also we have introduced a new parameter named “the effective height of the collector”, which presents the portion of the collector’s height illuminated by the reflector. This parameter is very important especially in case of PV solar fields, because it determines the domain of the concentrated solar energy over the surface of the PV panel.
Introduction
Many concentrator types are possible for increasing the flux of radiation on receivers.They can be reflectors or refractors.They can be cylindrical to focus on a "line" or circular to focus on a "point".Receivers can be concave, flat or convex.Tubular absorbers with diffuse back reflector, tubular absorbers with specular cusp reflector, plane receiver with plane reflectors; parabolic concentrator; fresnel reflector; array of heliostats with central receiver are types of concentrating collector configurations.The simplest and most inexpensive means for increasing the solar energy flux incident onto a surface is to attach one or more planar reflectors to the main harvester system.Concentration devices can produce elevated operating temperature under clear sky conditions, but require good optical components, more precise construction techniques and generally a mechanism for tracking the sun.A reflector augmenting a collector is, however, the best solution at utilizing both diffuse and beam (direct) radiation, while providing a moderate concentration with minimal tracking [1].The cost of the plane reflector is less than 5% of the cost of the PV system, while it can provide more than 15% yearly enhancement in solar energy collection inside soar energy devices.So, plane reflectors are very promising as an ideal method of improving the efficiency of PV modules with minor cost [2].The feasibility for the addition of flat reflectors to PV panels is techno-economically investigated in [3] for various applications (building attached PVs, ground installations, grid-connected or standalone units), various constructions and various PV types (mono-crystalline and amorphous silicon PV panels).External reflectors have been analysed theoretically and in some works also experimentally for a single solar collector in literature [4] [5] [6] [7] [8].Reflector in front of a collector is one of the options for improving the performance and cost-effectiveness of large collector fields.In solar fields the distance separating the rows has to be large to minimise shadowing effects early and late in the year.At high latitudes, a lot of solar radiation falling between the collector rows is not used in the summer.By introducing reflectors between the collector rows, most of this energy can be utilized by the collectors, reducing both the collector and land area requirements for a given load.Besides all these works the passivity side of reflector deployment in solar fields still under argumentation.
In this paper we present a modeling study of PV (photovoltaic) and/or thermal collectors with the aim of predicting the enhancement of the annual radiation harvested by a solar collector due to matching with the reflective surface, mounted diagonally between two adjacent rows.In this regards we have performed a theoretical analysis on a tilted collector and reflector system in order to determine the optimal angle of collectors corresponding to a specific solar field's design parameters and for any location on the Northern hemisphere.A cross section of a collector array with reflectors is shown in Figure 1.
We normally choose a tilt angle for solar collectors and plane reflectors in solar field installations, this increases the energy yield and decreases the losses due to collector and reflector dirtying compared to horizontal systems.Practically the reflector may be divided into two parts in such manner, that the upper part of the reflector ( , r upper
L
) is equal to the collector projection on the vertical plane ( , sin ), so the reflector can plait in order to allow technicians to pass through the space between the rows for cleaning or maintenance purposes.The shadow is not involved in this study because the reflector is not shadow causative in this arrangement, according to the study assumptions.
Solar Radiation on a Single Tilted Surface
For purposes of solar process design and performance calculations, it is often necessary to calculate the hourly radiation on a tilted surface.There are many models and softwares developed and utilized by researchers to estimate solar irradiance and PV-array performance in solar fields.The National Renewable Energy Laboratory System Advisor Model (SAM) is used four radiation models included in the TRNSYS radiation processor within SAM: Perez, Hay & Davies, Reindl, and isotropic sky.While in the EnergyPro is used the Reindl model.The model considers the anisotropy diffuse sky model formulated by Hay and Davis [9].This model is suitable for clear conditions, and most of the diffuse will be assumed to be forward scattered [10].It includes components of beam directly from the sun and diffuse irradiation from the circumsolar and the sky dome, and beam and diffuse irradiation reflected from the ground.The total solar radiation ( , t s T I ) on a tilted surface at slope ( s S ) from the horizontal for an hour as the sum of three components is given as: where
I
are the hourly diffuse radiation parts of the circumsolar and the isotropic on a horizontal surface, so the total diffuse radiation on a horizontal surface will be equal to the sum of these two components where sc G is the solar constant (1367 W⁄m 2 ), n is denotes to day of the year, and t z θ is the solar zenith angle at the time and day of interest.The 1 sin cos cos sin cos cos sin where: L denotes the local latitude, angle δ is the declination angle, and h is the hour angle: ( ) in where s t presents the solar time.And g is the ground-reflectivity, s sky F − and s g F − are the collector-sky and collector-ground view factor, respectively.For a single tilted surface, the view factors are [10]: In this paper, the ASHRAE clear-sky model is adopted to estimate the hourly beam normal ( t bn I ) and diffuse ( t d I ) solar radiation.The ASHRAE clear-sky model appears to be general enough for the objective of the paper; furthermore, we don't need to any information about the location of interest, except for the latitude angle.
The direct beam radiation and sky diffuse are calculated from the following formula [11]: where A, B and C are constants for every day and are given in Table 1 for the 21 st day of each month [11].
The optimum tilt angle had been early calculated in [12] and was found as: The annual solar radiation incident on a single tilted surface has been calcu- lated for comparison purpose.Figure 2 illustrates the optimum surface tilted angle and the corresponding annual solar radiation for various latitudes on the Northern hemisphere.
Geometry of the System
The flat concentrator system to which we will refer, to quantify the amount of radiation incident on the collector is extended diagonally from the top of the preceding row to the bottom of the succeeding row as it illustrated in Figure 1.
A cross-section of one row is depicted schematically in Figure 3, in where all the dimensions which we need for the simulation are presented for the collector-reflector system.
The values of the reflector height to the collector height ratio reflector tilted angle ( r S ) are calculated according to the field design parameters: the collector tilted angle, and distance separating the rows to the collector height ratio, respectively.According to Figure 3, the reflector tilt angle is given by the following formula: And, the reflector to collector height ratio is calculated from the following formula: The variation of the reflector parameters is calculated and plotted in Figure 4.
Modeling of Solar Radiation in Solar Fields Augmented by Plane Reflectors
Principally, the approach discussed herewith is based on the following assumptions: 1) The analysis is 2-D.This means no irradiance from the reflector is falling away from the sides of the collector, this assumption is acceptable for multi-rows large solar fields with very long length of the collector and reflector; 2) The sky-diffuse irradiation is assumed to be anisotropic, and is determined by using view factors presented in a previous work [13], for both the collector and the reflector; 3) The reflector surface is considered fully polished and therefore, the incident angle and reflect angle of the sunray are the same; 4) The collector is able to see only the sky and the reflector surface; 5) When the effective length ratio is greater than 5 the reflected rays considered parallel to the collector plane; 6) The collector and reflector surfaces are always illuminating and there is no shadowing effect; 7) The reflectivity of the reflector surface is constant and independent on the solar incidence angle.
The model considers the anisotropic diffuse sky model, which includes components of beam, diffuse irradiation, and beam and diffuse irradiation reflected from the reflector, as it depicted in Figure 5.The collectors are south facing ( angles and the diffuse components by the corresponding view factors, the summation of them will be the total solar radiation incident on the collector.Now our aim is to calculate the following: 1) The solar incident angles; 2) The view factors; and 3) The beam and sky-diffuse solar radiation components.
The Solar Incident Angles
There are only four independent variables for the solar field; the first two are the solar field design parameters , , the third is the location (latitude), and the fourth independent variable is, of course, the time represented in solar angles.
From our point of view, the incident solar angle represents the connection between the solar field design parameter and the position of the sun in the sky, and it will be a good indicator for analysing the problem, for this reason we choose the incident solar angle instead of solar altitude angle which adopted by others.
There are three incidence angles in this arrangement, those are: 1) Incident angle from the sun to the reflector ( , 2) Incident angle from the reflector onto the collector ( , 3) Incident angle from the sun to the collector ( , The solar incident angles , The incident angles have been calculated for Brack-Libya (L = 27.53˚N) at solar-noon for spring equinox, and for summer and winter solstices as a function of the solar-field design parameters and plotted in contour manner in Figure 6, Figure 7, and Figure 8.As presented in the Figure 6, the incident angle of the θ → ≥ ) the collector interested by the beam and diffuse components from the sun and the sky direct to the collector and the isotropic sky-diffuse reflected from the reflector only, because there are neither beam nor circumsolar incident on the reflector surface.Consequently the effective height ratio is set to be zero.
View Factors Calculations
Sky view factors for both the collector and the reflector are calculating according to [13].The equations are rearranged in terms of dimensionless parameters.The collector-sky view factor c sky F − and the reflector-sky view factor r sky F − are rewritten in the following forms: Using the view factor algebra the collector-reflector view factor is:
Dynamic Analysis of the Collector and the Reflector
Due to the continued change of the sun position on the sky, therefore, the solar radiation geometry is also changing.In this manner, we introduce five possible cases for the situation of the reflected beam irradiance from the reflector to the collector according to the concentration level.Accordingly, we illustrate the five cases graphically in Figure 9.In where: a) Zero concentration: when the reflector is blocked by the shadow of the previous row.b) Nonhomogeneous high-concentration: when all beam radiation incident on the reflector is reflected to a portion of the collector.c) Homogeneous normal-concentration: when all beam radiation incident on the reflector is reflected to the entire collector without any losses.d) Homogeneous low-concentration: when only a fraction of the beam radiation incident on the reflector is reflected to the entire collector with some radiation missing outside the collector.e) Zero concentration: when the reflected beam radiation from the reflector is parallel to the collector's plane.Figure 10 illustrates the dynamic analysis of all angles affecting the solar radiation incident directly from the sun on both surfaces of the collector and the reflector and reflected irradiance from the reflector to the collector.The diffuse-sky irradiances and its reflected component are independent of these angles.However, they are dependent on the view factors of the reflector-sky, collector-sky and collector-reflector view factors.
Knowing the solar incident angle on the reflector , t i r θ , which is a function of the time, day and the latitude, and according to Figure 5, one can determine the angles illustrated in Figure 10 from the following formulas: and 180 where t is the time, which presents the dynamic situation of the problem.
As regards the treatment of the problem we introduce a new quantity named "effective collector height ratio ".This quantity is an essential indicator for classifying the operation regime of the collector.This parameter indicates to the collector portion that is illuminated by the reflected irradiance from the reflector, and it is expressed as a dimensionless in the form: In this context, this parameter is considered as a tool to determine the solar energy situation of the collector-reflector system.Table 2 shows these situations according to the effective collector height ratio.Since all angles and view factors are determined, it is now possible to estimate the solar radiation incident on the reflector and the collector.
Solar Radiation Components Calculation
According to the above mentioned analysis, the definition of all irradiances ( , t c T I ) that strike the collector's surface, may be expressed as: where: the geometric factors , and r is the reflectivity of the reflector.
In cases (A and E) the ratio is set to be zero, which means the collector has not be received reflected beam and circumsolar irradiation and it has be received only isotropic diffuse from the reflector.
Results and Discussion
An MsExcel sheet has been prepared in order to calculate all variables involving in the simulation process.Table 3 presents the input variables and their values.The results have been obtained for a site (Brack -Libya) locates on latitude angle of 27.53˚N and longitude angle of 14.28˚E.
Applying Equation (17) to calculate the total solar radiation incident on the collector, the results obtained were plotted in Figure 12.The solar field design parameter , have been varied in order to see the effect of these parameters on the performance of the solar field that employed reflectors.Figure 12 shows the theoretical predictions of the hourly variations of the total solar radiation incident on the solar collector employed plane reflector for Brack El-Shati site at spring equinox and the solstices with distance ratio θ ) is large (see Figure 7) and Figure 13.Percentage of solar energy collection improvement with respect to a single collector tilted to optimum tilt angle which equal to 30˚ according to [12], for Brack El-Shati-Libya.
that occurred when the sun on the southern hemisphere during the winter solstice and also during the spring and fall equinoxes the opportunity of high concentration is large case (A, B and even D) and the effective height ratio is less than unity.While-in contrast-during the summer solstice, the sun is high in the sky dome and ( , t i r θ ) is small, the effective height ratio is relatively bigger and the re- flected energy tended to fail outside of the collector's area (cases C and E), as it evident from Figure 11.
The annual enhancement percentage in solar energy collection with respect to a single surface tilted angle 30˚ is depicted in Figure 13 for Brack El-Shati site.It is evidence from the Figure 13 that the optimum tilted angle of the collectors is increased with increasing the distance ratio; this relation can be established as a polynomial of fourth order with 2 1.0 R = by using MSExcel program: The obtained results show that, the annual solar energy collection increases by 72% for solar collector tilt angle of 70˚ and decreases by 4% for collector tilt angle of 10˚, when the distance separating the rows ratio
Conclusions
From our knowledge the solar radiation in stationary flat-solar fields is less than 5% than that in a single solar collector.This reduction in solar radiation is mainly due to the variation of the view factors of the solar field with respect to a single solar collector.A further 5% reduction in solar energy will occur due to the shadowing.A significant augmentation of solar radiation received by the collector surface reached to 75% by using plane reflectors with a distance separating ratio 2.0 The results were obtained by this work significantly coincident with previous experimental results carried by others (such as [5] [7]), which makes us confident in recommending the use of the proposed approach for the design and optimization purposes of solar fields.
Recommendations
Further investigation must be followed to examine the effects of the inhomogeneous solar radiation distribution through the solar collector on the electrical performance of the PV solar fields.As it is shown in this work, the operation regime of the solar collector is almost located in the case "A" where PV panels will suffer from the variation of solar radiation distribution.An economic study is recommended in order to determine the benefits of utilization of plane reflectors in solar fields of PV panels.Furthermore, the assumption that the output is proportional to irradiation will probably overestimate the output from PV modules with reflectors, since increased irradiation from the reflector increases the module
Figure 1 .
Figure 1.Deployment of stationary solar collector arrays with plane reflectors on a solar field.
R
is a geometric factor which presents the ratio of beam radiation on the tilted surface to that on a horizontal surface at any time,
Figure 2 .
Figure 2. The optimum surface tilted angle and the corresponding annual solar radiation vs. latitude angle.
Figure 3 .
Figure 3. Side view of one collector-reflector row, where: 1 cos c c L L S = ,
Figure 4 .
Figure 4.The values of the reflector design parameters (a): ( ) r S and (b): 0 c ψ = ) and the reflectors are north facing ( 180 r ψ = ).The configuration is considered to be of infinite length and with variable solar field design parameters , the beam components by the corresponding incident
Figure 5 .
Figure 5. Solar radiation components and the three solar incident angles.
S
) for the reflector and the collector respectively, at the time and the location of interest.Accordingly ( ,t i r c θ → ) is cal- culated from Figure10 below:
Figure 6 .
Figure 6.The solar incident angle on the collector surface ( , t i c θ ).
Figure 7 .
Figure 7.The solar incident angle on the reflector surface ( , t i r θ ).
Figure 8 .
Figure 8.The reflected solar incident angle from the reflector on the collector surface ( , t i r c θ → ).
angle c S .While, the incident angle of the reflector , , and so the incident angle from the reflector onto the collector ( ,t i r c θ → ).All incident angles range 0 ,90 , in case of the reflected incident angle be-
Figure 9 .
Figure 9.The five basic situations of incidence beam solar radiation onto the reflector and the collector (red arrows), and reflected from the reflector onto the collector (blue arrows).
Figure 11 2 .Figure 11 .t θ accordingly the value of 3 tθ
Figure 11 presents the variation of the effective height ratio 12:00 , c eff c . The behaviour of the solar incident angle ( the reflector presents the key for understanding the process.When ( , t i r
Figure 12 .
Figure 12.Hourly total solar radiation incident on a solar collector employed a flat reflector for spring equinox and for summer and winter solstices for various solar field design parameters.
equal to 1 . 5 .
The maximum increase in solar energy collection was 32.7% for solar collector tilt angle of 60˚ and decreases about 8.6% for solar collector tilt angle of 90˚.The situation was dramatic for in the performance of the collector.In the synoptic Figure14the improvement of solar radiation due to the plane reflector in solar field are depicted as a function of the solar field design parameters ( C S and c X L ) and for many latitudes in the Northern hemisphere.The importance of the Figure 14 is, one able to determine the optimum solar collector
Figure 14 .
Figure 14.The obtained improvement of solar radiation vs. the solar field design parameters ( C S and
Figure 15
Figure 15.Z-axis presents the optimum collector tilted angle are playing crucial role in the augmentation process, in addition, the results show that the percentage improvement in solar radiation is independent on the latitude angle and the optimum solar collector tilt angle is-almost-for all latitudes are the same.Nonhomogeneous solar radiation distribution through the collector issues appears during the day which is presented a serious problem especially in the case of the PV panel solar field. m
Table 1 .
Constants for ASHRAE equations for the 21 st day of each month.
Table 3 .
The input parameters and variables included in the MsExcel sheet. | 5,210.8 | 2017-07-06T00:00:00.000 | [
"Engineering",
"Physics",
"Environmental Science"
] |
Two New Extended PR Conjugate Gradient Methods for Solving Nonlinear Minimization Problems
In this paper, we have discussed and investigated two nonlinear extended PR-CG method which use function and gradient values. The two new methods involve the standard CG-methods and have the sufficient descent and globally convergence properties under certain conditions. We have got some important numerical results by comparing the new method with Wu and Chen PRCG-(2010) method in this field.
Introduction
This paper considers the calculation of a local minimizer x* say, for the problem: is a smooth nonlinear function (of n variables) and its gradient vector is available are calculated but the Hessian matrix is not available. At the current iterative point k x , the Conjugate Gradient (CG) method has the following form: where k is a step-length; k d is a search direction; k is a parameter.
Standard algorithms for solving this problem include CG-algorithms which are iterative algorithms and generate a sequence of approximations to minimize a function ) ( F x and their very low memory requirements. However, this paper considers a more general model than the usual quadratic function ) ( ) (5) and (6) are called the "Standard Wolfe" and "Strong Wolfe" conditions, respectively. When ,quadratic functions and exact line searches are used, all the above formulas in (4) are equivalent. However, these formulas very according to general functions. For general functions, [22] proved the global convergence of PR method with exact line search. On the other hand, the PR and HS methods perform similarly in terms of theoretical property. Nevertheless, [16] showed that the PR and the HS methods can cycle infinitely without approaching a solution, which implies that they do not have globally convergence.
In this paper, we have proposed, two new special formulas keeps the property of PR method, namely, if a very small step is generated the next search direction tends to the Steepest Descent (SD) direction, preventing a sequence of tiny steps from happening. Furthermore, finite quadratic termination is retained for the new methods. Since the sufficient descent condition is a property of great importance for the global convergence analysis of any CG-method, we have modified the conjugacy parameter of [21] to implement the nonquadratic rational model which satisfies the sufficient descent property and the modified Wolfe-Powel conditions introduced by Andrei [6] we illustrate this condition in section 4 . In addition, the global convergence property of the new proposed CG-method is discussed and a set of numerical results presented show that the new proposed method is efficient.
Extended CG-Methods For Non-Quadratic Models.
Over years, various authors have published works in this area, In this paper, a more general model than quadratic one is suggested as a basis for a CGalgorithm. If q(x) is a quadratic function, then a function F(q(x)) is defined as a non-linear scaling of q(x) if the following invariancy condition holds: ( 7) where, * x is the minimizer of q(x) with respect to x for more details see [19] and f is monotonic increasing, may be better to represent the objective and thus it gives an advantage to method based on this model. In order to obtain better global rate of convergence for minimization methods when applied to more general functions than the quadratic.
The following properties for ) (x f are immediately derived from the above condition: • Every contour line of q(x) is a contour line of x is a minimizer of q(x) , then it is a minimizer of in at most n step has been described by Fried [11].
ERCG-Method (Al-Bayati,1993). [2]
Al-Bayati's, 1993 non-quadratic model is defined as the quotient of two quadratic functions and so belongs also to the class of rational functions , Al-Bayati's rational function model was considered by:
The Special Cases.
In this paper ,we introduced the two special cases of AL-Bayati's (1993) and Tassopoulos and Storey (1984) extended CG-method which are invariant to nonlinear scaling of quadratic rational functions are proposed. The first investigated model is defined as the quotient of two quadratic functions and so belongs also to the class of rational functions a special of AL-Bayati's rational function model was considered by: ….. (11) Two New Extended PR Conjugate Gradient Methods for Solving Nonlinear … 76 From (7) We can rewritten (11) as: Where the function f is defined as nonlinear scaling of ) (x q and the invariancy property to nonlinear scaling (7) holds and ) (x q is defined in (9) is the quadratic function then it determines the solutions min x in a finite number of iterations not exceeding (n). It is shown the one-dimensional problem ( ) and k d is a search direction that the following updating process by the Boland theorem [10] to convert the quadratic model to a non-quadratic model in (12) we can write: (13) Since from the (12c) we have: If substituting (13) in (14) we get: Similarly the special case of the rational function of Tassopoulos and Storey [20] in (8a). In this paper we shall state the rational function of AL-Assady and Shakory [18] considered by: Since from the (12) we have: If substituting (16) in (17) we get: The (18) and (15) are called the special case of rational function model of Tassopoulos and Storey (1984) and AL-Bayati's (1993) ,respectively.
Two New Combined of Rational Functions.
A. We introduce the combined of two rational function as a convex combination of the special case of Tassopoulos and Storey in the equation (18) and the rational function AL-Bayati's (1993) in the equation (10) nonquadratic model to be investigated here is considered by :
B.
We introduce another combined of two rational function as a convex combination of the special case of AL-Bayati's in (15) and the and AL-Bayati's in (10) non-quadratic model to be investigated here is considered by :
Wu and Chen (2010) CG-Method.
In this section, we are going to present the recent work of the two wellknown Scientist Wu and Chen in (2010). They introduced several wellknown CG-formulas. The conjugacy parameters of these CG-methods are given by; , respectively by making use of the Powell's restarting criterion and the Armijo-type line search defined by: (22) They proved that all the above CG-methods satisfy the sufficient descent condition and have the global convergence property for more details see [21].
A New Extended CG-Method. 4.1 Transform Quadratic model to Non-quadratic.
Consider the following quadratic model we proceed as in [21]: (19)- (20) to get two new extended CG method whose conjugacy parameters are defined by:
An Acceleration Scheme of the Line Search Parameter.
In the CG-methods the search directions tend to be poorly scaled and as a consequence the line search must perform more function evaluations in order to obtain a suitable step-length k . In order to improve the performances of the CG-methods the efforts were directed to design procedures for direction computation based on the second order information. Jorge Nocedal [14] pointed out that in CG methods the step lengths may differ from 1 in a very unpredictable manner. They can be larger or smaller than 1 depending on how the problem is scaled. Numerical comparisons between CG methods and the limited memory QN method, by Liu and Nocedal [13], show that the latter is more successful [8]. Here, we have pointed out Andrei's [7] acceleration scheme; basically, this modifies the step length in a multiplicative manner to improve the reduction of the function values along the iterations [5,6].
Outline of The Two New Extended CG-Method.
Step 1: Given
) is an index of the algorithm
Step 2: Set k=1; Step 3: Using the modification WP line search conditions which fully described by Andrei (2009) determine the step length k , such that, compute: Acceleration scheme, compute , , then compute Step 4: Compute Step 5: If Powell restarting, , satisfied then set: (27)), go to Step 2.
2 Theoretical Properties for the Two New Extended CG-Method.
In this section, we focus on the convergence behavior on the 2 1 , New k New k methods with inexact line searches. Hence, we make the following basic assumptions on the objective function.
Assumption.[21]
f is bounded below in the level set U of the level set 0 x L , f is continuously differentiable and its gradient f is Lipschitz continuous in the level set 0 x L , namely, there exists a constant L> 0 such that: for all x, y The third term from the last equation can be simplified as: We now prove the theorem by contradiction and assume that there exists some constants > 0 such that For initial direction we have: Since our function f is uniformly convex function either in the quadratic or in the non-quadratic regions, then there exist a Lipschitz constant L >0 and a constant, 0 such that: For inexact line search using Wolfe-Powell conditions (5) and (6) we have: Powell and criteria which are defined as: From [17] Powell restarting criterion (45) we have: Using (45) Multiplying the search direction of (27) by T k g yields: For inexact line search using Wolfe-Powell conditions (5) and (6) we have: Thus our new proposed extended CG-method has sufficient descent directions using inexact line searches under the condition that Powell restarting condition must be used.
Theorem
Suppose that Assumption 4.3 hold. Consider the method (2)-(3) with the following three properties: , for more details see [12]. Therefore, the method has a global convergent property by satisfying the conditions of Zoutendijk theorem and the line search satisfy the strong Wolfe condition then from Gilbert and Nocedal in [12] these method is global convergent. | 2,259.6 | 2014-03-01T00:00:00.000 | [
"Mathematics"
] |
EFFICIENCY OF DATA MINING TECHNIQUES FOR PREDICTING KIDNEY DISEASE
- Chronic kidney disease is an aging problem in the current growing population. Kidney disease surveillance and prediction is very important for patients to provide adequate and appropriate treatment at the right time. Data mining can extract interesting patterns for gigantic medical databases. Patients with kidney disease can be automatically analyzed from their disease data taking into account prior predictions. Though medical data is heterogeneous in nature including text, graphics and images, unwanted data can be removed to provide useful medical information on a patient. Medical data mining can detect disease patterns and predict severity of a patient's disease. Conformist theories are more pertinent than probabilistic theories for results as precise results and inferences become a necessity to save a patient’s life. Fuzzy systems are generally used as they produce results based on mathematics, instead of probabilistic arbitrations like neural networks. The paper proposes new algorithm Improved Hybrid Fuzzy C-Means (IHFCM) which is an improvisation of FCM with Euclidean distances to predict kidney diseases in patients.
Basma Boukenze et.al [2] pre-processed data with conversions and data mining methods to gain knowledge about the interaction between measurement parameters and the survival of a patient. Two data mining algorithms were used to form decision rules in extracting knowledge and predict the survival of patients. They explained the significance of exploring important parameters using data mining. Their new concept was implemented and tested using dialysis data collected from four different sites. Their method also reduced the cost and effort in selecting patients for clinical trials. The patients were selected based on predicted results and significant parameters found in their analysis. Neha Sharma et al [3], detected and predicted kidney diseases as a prelude to proper treatment to patients. The system was used for detection in patients with kidney disease and the results of their IF-THEN rules predicted the presence of a disease. Their technique used two fuzzy systems and a neural network called a neural blur system, based on the result of the input data set obtained. Their system was a combination of fuzzy systems that produced results using accurate mathematical calculations, instead of probabilistic based classifications. Generally results based on mathematics tend to have higher accuracies. Their work was able to obtain useful data along with optimizations in results. Veenita Kunwar et al [4]. In their study predicted chronic kidney disease (CKD) using naive Bayesian classification and artificial neural network (ANN). Their results showed that naive Bayesian produced accurate results than artificial neural networks. It was also observed that classification algorithms were widely used for investigation and identification of CKDs. Swathi Baby P et al [5] demonstrated that data mining methods could be effectively used in medical applications. Their study collected data from patients affected with kidney diseases. The results showed data mining's applicability in a variety of medical applications. K-means (KM) algorithm can determine number of clusters in large data sets. Their study analyzed tree AD, J48, star K, Bayesian sensible, random forest and treebased ADT naive Bayesian on J48 Kidney Disease Data Se and noted that the techniques provide statistical analysis on the use of algorithms to predict kidney diseases in patients.
III. PROBLEM FORMULATION Probability theory cannot be used to obtain the results in prediction of kidney diseases as it involves the patient's life and the exact results are a necessity. Statistical methods, Bayesian classification or association rule based predictions cannot be used to predict CKD as the results obtained may be less accurate. Predicting disease can save a patient's life and if detected early can help proper cure of the disease. Thus a need to evolve CKD prediction with new techniques.
IV. PROPOSED WORK Diseased kidneys are increasing in an aging population making it imperative to monitoring or prediction diseased kidneys. General predictions are based on a set of if then rules on kidney datasets. Erroneous predictions of CKD can lead to loss of life. The proposed a new technique IHFCM is used for predicting and detecting kidney disease in a patient data set.
V. METHODOLOGY A. Fuzzy Model
Fuzzy grouping is based on generation of graphs for each pattern within the group. Fuzzy modeling can match human reasoning models and manage data. The main advantages of fuzzy logic include its simplicity and flexibility. Fuzzy logic can handle inaccurate and incomplete data where traditional statistical models may fail. A fuzzy system can be any model of a complex nonlinear function and provides transparency with explanation on rules. These rules can be potential clinical guidelines.
B. Fuzzy C Means
The fuzzy c-means (FCM) algorithm is a traditional and classical image segmentation algorithm. It is a method that allows clustering, where data may belong to two or more clusters The FCM algorithm focuses on minimizing the value of an objective function that measures the quality of the partitioning a dataset into clusters. It produces an optimal partition by minimizing the weights within a group sum of squared error objective function. It is frequently used in pattern recognition. The fuzzy C-means algorithm is listed below in
D. Proposed Improved Hybrid Fuzzy C Means Clustering Algorithm (IHFCM)
The fuzzy c-means is introduced by Ruspini and then extended by Dunn and Bezdek and is widely used as clustering analysis, pattern recognition and image processing in Fuzzy C Means Clustering Algorithm (FCM). It is based on the K-means and the basic idea of FCM that each data point belongs to the membership in the degree of poor clustering, and K means that each data point belongs to a particular group or not. So FCM uses fuzzy partitioning so that when you can belong to multiple groups, the members are between 0 and 1. However, through the degree of data provided by the degree of membership, FCM still uses the cost function to try to split the data set. When minimized. It makes the matrix member having a U element value between 0 and 1. The algorithm works iteratively through the preceding two conditions until the no more improvement is noticed. In a batch mode operation, FCM determines the cluster centers i, c and the membership matrix U using the following steps: Input: Feature extracted CT scan kidney segmented image Output: given image has Kidney Disease or not kidney disease Step 1: Set the number of clusters Step 2: Set the Fuzzification parameter, image size and ending condition.
Step 3: Initialize randomly the fuzzy cluster and conditions.
Step 4: Set the loop condition initialize by 0 Step 5: Calculate the weighted fuzzy factor using Euclidean distance measure.
Step 6: Modify the segmented matrix M= {M ij } using Euclidean Distance (d) Step 7: Modify the Cluster conditions using fuzzy membership function (MF) Step 8: If (MAX|MF new -MF old | < End Condition) then Stop Step 9: otherwise increment Loop condition +1 and go to step 5. Where MF= [MF 1 , MF2… MF C ] are membership function of cluster condition. At the end point, a defuzzification process takes place to convert the fuzzy image to crisp segmented image. IHFCM can be applied to Identifying a disease in a patient's dataset and even be used for Drug Activity Prediction.
VI. EXPERIMENTAL RESULTS This work is done on MATLAB which can manipulate matrices, product functions and data, implement algorithms, create user interfaces, and interact with programs written in other languages. The experimental IHFC is worked on MATLAB. The data set is extracted from the reference point UCI library machine. In the UCI machine learning library they are in the machine learning community used in the machine learning algorithm to conduct an empirical analysis of the field of database theory and data generation. The document was created by David Aha in the 1987 FTP file and other graduate students at the University of California, Irvine. Since then, it has been widely used by students, educators and researchers from major sources of data collection machines around the world.
A. Fuzzification Score
The algorithm calculates the fuzzy C meaning as the diffuse score for each value in the corresponding table of the contents of the query that is entered as a score. The higher the score, the more similar the string. A score of 1.0 or 0.9 means that the fuzzy score results in a highly risky clustering. 0.0% means that the corresponding symptoms have a risk level that is less affected or is not at risk. The user can enter the minimum and highest possible risk factors that are set to contact the doctor and the base, the individual gives each query score, FCM is divided into two categories with the lowest and highest levels found again with the result with the given range of values Find the minimum and maximum scores given to their limits. Thus, FCM can provide three low-risk scores for finding high-risk results with fuzzy scores, fuzzy average sub-risk and cluster-based results.
B. Results
The performance of FCM is evaluated by statistical measures like sensitivity, specificity and accuracy to illustrate the normal life style score. These metrics also enumerate how the test was good and consistent. Sensitivity evaluates the normal life style score correctly at detecting a disease positively. Specificity measures how the proportion of patients without disease can be correctly ruled out. The objective function of IHFCM is depicted in Figure 3. The comparative performance of the algorithms is listed in table 1 and
VII. CONCLUSION
The proposed IHFCM is an extension of FCM and is applied for locating kidney disorders in patient records. The paper demonstrates that correct adjustment to FCM can help build a new strategy for discovering unusual and traditional cases. Initial pre-processing of IHFCM is deleting duplicate records. Results of clustering which obtained from 300 patients showed that FCM based clustering algorithms achieve higher accuracy than most existing algorithms. The proposed IHFCM's performance has been proved clearly in terms of accuracy. | 2,298.8 | 2017-10-31T00:00:00.000 | [
"Computer Science"
] |
On the Performance of a Wireless Powered Communication System Using a Helping Relay
This paper studies the outage performance and system throughput of a bidirectional wireless information and power transfer system with a helping relay. The relay helps forwardwireless power from the access point (AP) to the user, and also the information from the user to the AP in the reverse direction. We assume that the relay uses time switching based energy harvesting protocol. The analytical results provide theoretical insights into the effect of various system parameters, such as time switching factor, source transmission rate, transmitting-power-to-noise ratio to system performance for both amplify-and-forward and decode-and-forward relaying protocols. The optimal time switching ratio is determined in each case to maximize the information throughput from the user to the AP subject to the energy harvesting and consumption balance constraints at both the relay and the user. All of the above analyses are confirmed by Monte-Carlo simulation.
Introduction
Recently, radio frequency (RF) signal based wireless energy transfer (WET) has emerged as a perpetual and costeffective solution to power wireless devices, such as mobile sensors, electronic tags, etc. [1].While numerous works have focused on WET systems to optimize the energy harvesting process at energy receivers [2], the authors of this paper are more interested in another line of WET research, where WET could be integrated with wireless communication by exploiting the dual use of RF signals.Especially, we focus on wireless powered communication (WPC) [3], where the energy for wireless communication at the device is obtained via the WET technology.This advanced technology has been deployed and investigated in various wireless system models, including cellular networks [4], relay systems [5], [6], cognitive radio networks [7], [8].
In last decade, wireless sensor networks have been more and more attracted by research community, due to their ability to carry out different kinds of tasks, from traffic monitoring, agriculture monitoring, to smart home and health-care applications.For these networks, network life time is a critical aspect to the success of the system.By some previous research, battery life time is the bottleneck in determining the life time of the whole system.In [9], three wireless powered sensor network models for infrastructure monitoring application have been proposed and their performances have been investigated.Another wireless powered sensor network model was presented in [10], in which a number of sensor nodes send common information to a far apart information access point via distributed beamforming, by using the wireless energy transferred from a set of nearby multi-antenna energy transmitters.
Both of the works in [9] and [10] only introduce normal sensor networks without the helping of relay nodes.Furthermore, the RF energy transmitters in those works are independent of the information transfer process.This would increase the cost of implementation of these models in practice.In [11], the authors have tried to overcome this drawback by considering a new WPC system, where a wireless user communicates with an access point (AP) assisted by a bidirectional relay.The user and the relay are both powered by the RF energy from the AP.Here, the role of the relay is to forward the energy from the AP to the user, as well as to forward information from the user to the AP.However, the authors in [11] only considered the case that channel gains are constant, and estimate the maximum achievable throughput of the system.Because of this limitation, there is a large difficulty to apply this result to practical sensor networks.In addition, the work in [11] only considers amplify-andforward as the relaying strategy.
Continuing to the work of [11], in this paper we provide a rigorous analysis on the same wireless powered sensor network model.We apply a Rayleigh distribution model for the channel gains between nodes, including the AP, relay node and the wireless user.For information transfer, both amplify-and-forward (AF) and decode-and-forward (DF) relaying protocols are investigated.Regarding to the energy harvesting protocol, we focus on time switching (TS) strategy at the relay.The outage probability and the average throughput of the system are derived mathematically.The optimal time switching factor to maximize the system throughput is obtained via numerical algorithm.To verify the analysis mentioned above, Monte-Carlo simulations are also conducted and the results are reported in this paper, too.
The rest of this paper is organized as follows.The next section introduces the system model that we are going to analyze.The detailed performance analysis is provided in Sec. 3. The numerical results to support the analysis are given in Sec. 4. Finally, Sec. 5 concludes the paper.
System Model
We consider a wireless powered system as illustrated in Fig. 1, where a mobile user is intended to send information to the AP with the assistance of a relay R. Assume that both the user and relay R have no other energy supply but solely the energy harvested from the AP.Furthermore, we assume that direct connection between the AP and the user is so weak, hence, the only available communication path as well as power transfer path is via the relay R. The relay serves the dual roles of both energy relaying from the AP to the user and information forwarding from the user to the AP [11].To initialize the communication process, a sufficient amount of initial energy is stored in the battery to conduct the first transmission block before energy harvesting, as in [12].After that, the energy consumed by the user/relay is kept lower than or at most equal to the harvested energy amount during each block, thus no further manual battery replacement/recharging is needed.
All nodes are assumed to operate in half-duplex mode, and either amplify-and-forward (AF) or decode-and-forward (DF) relaying strategy can be used at the relay for information transferring.Regarding to the channel model, we consider the case that perfect channel state information (CSI) is available at the relay and the AP.Let h and g denote the channels from the AP to the relay and from the user to the relay, respectively.In addition, we assume for simplicity that these channels are reciprocal.Different from the work in [11], all channels here experience Rayleigh fading and keep constant during each transmission block so that they can be considered as slow fading.As a result, |h| 2 and |g| 2 are an exponential random variables with parameters λ h and λ g , respectively.
For energy harvesting, we employ the time switching relaying (TSR) protocol, which is more convenient to implement in practice.As shown in Fig. 2, the total symbol duration T is divided into three intervals with the lengths of αT, (1 − α)T/2, and (1 − α)T/2, respectively, where 0 < α < 1 denotes the time-switching ratio.The first interval corresponds to the energy harvesting phase at the relay R, in which the AP wirelessly sends its energy to R with power P ap .Then, the total energy harvested at R during each block is given by E r = ηP ap .|h| 2 .αT,where 0 ≤ η ≤ 1 is the energy conversion efficiency.The second phase of duration (1 − α)T/2 corresponds to the information transmission from the user to the relay.In the third phase of the transmission block, R forwards an amplified or decoded signal to the AP and also forwards energy to the user.We assume that the circuit power consumption is negligible as compared to the radiation power, which is reasonable for low-power devices such as sensor nodes.
Performance Analysis
In this section, the throughput and outage performance of the proposed system are analyzed mathematically.The impact of time-switching factor on system performance is investigated.We consider both AF and DF protocols in our analysis.
Amplify-and-Forward Protocol
Let x u denote the transmitted signal from the user during the second phase and P u denote the power of this signal.The received signal at R during this phase is espressed as where n r ∼ N (0, N 0 ) denotes the Gaussian distributed noise at the relay R.
During the third phase of transmission, the relay amplifies the received signal from mobile user and forwards it to both the AP and the user.While the AP receives this signal for the purpose of getting information message, the user receives the same signal for energy harvesting purpose.The received signal at the AP during this phase is written as where x r is the signal transmitted by the relay, which has the power of P r , and n d is the zero-mean Gaussian noise at the AP with variance N 0 .Because the transmit power of the relay comes from the energy supplied by the AP in the first phase, we must have [11] where k = 2α 1−α .The signal transmitted by the relay is an amplified version of y r : x r = βy r .
According to energy conservation law, the energy consumed by the relay cannot exceed its available energy, which yields [11] β = 1 Now, we can substitute (1), (4), and ( 5) into (2) and get From ( 6), the signal-to-noise ratio at the AP can be computed by Let's move on to determine P u .We know that the received signal at the mobile user during the third transmission phase is y u = gx r = g √ βy r .Hence, the energy harvested during this phase can be determined by So, the transmit power of the mobile user during the second phase is expressed as By substituting ( 9) into (7) and doing some algebra, we obtain the overall SNR for AF protocol: Assume that the source transmits at a constant rate R, then γ = 2 R −1 is the lower threshold for SNR.Here, the outage probability P out and the average throughput of the system can be evaluated by [6] The main contribution of this paper is to derive the closed-form expression of the outage probability and average throughput of the system of interest, as well as to figure out the optimal time-switching factor for energy harvesting.The results for AF protocol are formally stated in the following theorems.Theorem 1 provides the exact integral forms for the outage probability and throughput of the proposed system with AF protocol.In Theorem 2, closed-form approximations of the outage probability and throughput in terms of Meijer function are derived for high source-power-to-noiseratio regime.
Theorem 1 (AF Protocol) For the AF protocol, the outage probability and the average throughput of the proposed system can be expressed as
and where δ = kη
Theorem 2 (AF Protocol -Closed-form approximation)
At high P ap /N 0 regime, the outage probability and average throughput of the proposed system with AF protocol can be respectively approximated to and where G m,n p,q (•| • • • ) is the Meijer function (Sec.9.3 of [13]).Proof 2 See Appendix B.
Decode-and-Forward Protocol
For DF relaying protocol, the data communication is divided into two separating hops, which do not depend on each other.Hence, the outage occurs if and only if either the source-relay path or the relay-destination path fails to satisfy the corresponding SNR constraint.Different from the AF protocol, the message transmitted by the relay during the third transmission phase is the decoded message xr , instead of x r , and the transmit power of the relay in this phase is the same as the one given in (3).Hence, the energy harvested by the mobile user during the same transmission phase is As a result, the transmit power of the mobile user in the second phase is the same as in (9).
According to the equations ( 1) and ( 2), the SNR values at the relay R and the AP are respectively determined by The outage probability of the system can be written as Now we can claim the following theorem on the outage probability and the average throughput of the system of interest.
Theorem 3 (DF Protocol) For the DF protocol, the outage probability and the average throughput of the proposed system can be expressed as
and where Γ(α, x; b, β) complete gamma function, which is defined in [14], and x 0 , y 0 are defined by Proof 3 See Appendix C.
Optimal Time-Switching
To find the optimal time-switching factors that give the best performance in terms of outage probability or average throughput, we solve the equations dP out (α) dα = 0 and dR(α) dα = 0, respectively, where P out (α) and R(α) are outage probability and throughput functions with respect to the time-switching factor.
By investigating the outage probability functions with respect to α for both AF and DF, we can easily see that these are non-increasing functions.That means, the best outage performance is obtained when we exploit energy harvesting at full-scale.However, we should keep in mind that this outage performance only based on the comparison of power between signal and noise.It ignores other factors of communication process.In practice, we cannot set α to 1 because it means that no communication data is transferred.
Hence, the average throughput should be a more reasonable performance factor to be optimized.By plotting the throughput functions for AF and DF protocols versus α, we learn that these functions are concave functions, which have a unique maxima on the interval [0, 1].The optimal factor α * can be found numerically by some iterative methods, for instance, Golden section search method [15].
Numerical Results and Discussion
In this section, we conduct Monte Carlo simulation to verify the analysis developed in the previous section.For simplicity, in our simulation model, we assume that the sourcerelay and relay-destination distances are both normalized to unit value.Other simulation parameters are listed in Tab. 1.
Amplify-and-Forward Protocol
In Figures 3 and 4, the achievable throughput and outage probability of the system with AF protocol are plotted against P ap /N 0 ratio with the data rate set to be 3 bps.The time-switching factor α is chosen to be 0.3 and 0.7.It's can be observed that the outage probability is a decreasing function with respect to P ap /N 0 , while the throughput grows with P ap /N 0 .In addition, the simulation and the analysis curves are overlapping.The approximate outage probability and throughput are also plotted in these figures.They are close to the exact curves, especially when P ap /N 0 is large.This confirms the correctness of our analysis in the previous section.The impact of time switching factor on system performance with AF protocol is illustrated in Fig. 5 and 6.In this experiment, P ap /N 0 is set to 5 dB, and the rate can be varied at 3 bps, 2 bps, and 1 bps.It can be observed that the outage probability is reduced when we increase the value of α.On the other hand, the simulation result shows that there exists a unique time switching factor at which the average throughput is maximized.Indeed, this optimal factor can be found iteratively using numerical methods.
Decode-and-Forward Protocol
For decode-and-forward protocol, we also have similar results about the impact of various parameters, such as P ap /N 0 and α on the average throughput and the outage probability of the system.Specifically, Fig. 7 and 8 respectively plot the outage probability and throughput against P ap /N 0 , while Fig. 9 and 10 show the dependence of these performance characteristics on time-switching factor α. than AF protocol in terms of both outage probability and throughput, because the noise at relay is eliminated in DF protocol, while it's accumulated and amplified in AF protocol.
Optimal Time-Switching Factor
Finally, the optimal values of α at different values of source-power-to-noise-ratio for both AF and DF protocols are shown in Fig. 13.We can see that the α value that optimizes the throughput has tendency to decrease when P ap /N 0 increases.
Conclusions
In this paper, we investigate the performance of a new WPC system with a bidirectional information/energy forwarding relay in the Rayleigh fading environment.Two relaying protocols based on AF and DF strategies at the relay are considered in our work.For practical orientation, we employ the time switching protocol for energy harvesting.The exact-forms of outage probability as well as the average throughput of the proposed system are derived rigorously.Numerical results are provided to verify our analysis.The results show that the outage probability decreases as the time switching factor increases, while there is a unique value of time switching factor such that the throughput is maximized.For comparison between relaying protocols, the DF protocol is slightly better than its counterpart.
While the motivation of this paper comes from the energy problem in wireless sensor networks, the analysis obtained in this work does not limit to wireless sensor networks themselves, but can be applied for a wide range of wireless applications that employ the relay-node idea.Due to this reason, some specific issues related to sensor networks have not been considered in this paper.For example, the energy required by sensor nodes when collecting data or making measurement should be taken into account.In that case, the energy source not only comes from the information source node but also from other available nodes.This harvested energy can be modeled as a randomly varying variable.That should be our further work in this topic.In addition, we can take into account other factors such as CSI error and hardware impairment.
From (19), the outage probability can be rewritten as By substituting this into (C.1) and using the definition of extended incomplete gamma function (formula (1.9) in [14]) we obtain (20).Finally, (21) is obtained by including (20) in the definitive formula of average throughput.
Fig. 3 .
Fig. 3. Outage probability versus source power to noise ratio for AF protocol.
Fig. 7 .
Fig. 7. Outage probability versus source power to noise ratio for DF protocol.
Fig. 8 .
Fig. 8. Throughput versus source power to noise ratio for DF protocol.
Figures 11 and 12
Figures 11 and 12 compare the performance of two protocols that are considered in this paper.The results show that the DF protocol is slightly better
Fig. 11 .
Fig. 11.Outage probability of AF and DF protocol at rate 3 bps. | 4,169.2 | 2017-09-15T00:00:00.000 | [
"Engineering",
"Computer Science"
] |
Low-cost single-pixel 3D imaging by using an LED array
We propose a method to perform color imaging with a single photodiode by using light structured illumination generated with a low-cost color LED array. The LED array is used to generate a sequence of color Hadamard patterns which are projected onto the object by a simple optical system while the photodiode records the light intensity. A field programmable gate array (FPGA) controls the LED panel allowing us to obtain high refresh rates up to 10 kHz. The system is extended to 3D imaging by simply adding a low number of photodiodes at different locations. The 3D shape of the object is obtained by using a noncalibrated photometric stereo technique. 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Eng. 19(1), 191139 (1980). 44. P. Favaro and T. Papadhimitri, “A closed-form solution to uncalibrated photometric stereo via diffuse maxima”, IEEE Conference on Computer Vision and Pattern Recognition, 821–828 (2012). 45. A. Yuille and D. Snow, “Shape and albedo from multiple images using integrability”, IEEE Conference on Computer Vision and Pattern Recognition, 158–164 (1997). 46. H. Hayakawa, “Photometric stereo under a light source with arbitrary motion,” J. Opt. Soc. Am. A 11(11), 3079– 3089 (1994). 47. P. N. Belhumeur, D. J. Kriegman, and A. L. Yuille, “The bas-relief ambiguity,” Int. J. Comput. Vis. 35(1), 33– 44 (1999). 48. http://www.cvg.unibe.ch/tpapadhimitri/ Vol. 26, No. 12 | 11 Jun 2018 | OPTICS EXPRESS 15624
Introduction
Digital cameras based on CCDs or CMOS sensors are key devices in most of the current imaging techniques. However, light sensors with a pixelated structure are not always necessary. In fact, imaging methods requiring raster scanning techniques, such as confocal or two-photon microscopy, use in general a single high sensitivity photosensor. Some other imaging techniques can benefit from light detectors with no spatial structure. This is the case of imaging with low light levels or with electromagnetic radiation in spectral ranges out of the visible region, particularly in the near IR and terahertz spectral bands, where it is more difficult to make a pixelated sensor [1][2][3]. Imaging with single-pixel detectors can be useful also when the purpose is to measure the spatial distribution of several optical parameters of the light beam simultaneously, not only the intensity but also the polarization state [4], the phase [5,6], the color [7][8][9][10], or the spectral content [11][12][13][14][15][16]. In particular, single-pixel recording strategies are well adapted to provide depth information and have been applied with success in many different 3D imaging techniques. Single-pixel imaging techniques [25,26] are based on illuminating the scene with a sequence of light structured patterns while the light reflected or transmitted by the objects is recorded by a single photosensor such as a photodiode or photomultiplier tube (PMT). Alternatively to this active imaging system, where a spatial light modulator (SLM) modifies the input light beam, in the passive configuration the light coming from the object is sampled sequentially by a set of masks codified in the SLM. The image is then computed numerically by using different mathematical algorithms such as a simple linear superposition, a change of basis transformation, or a correlation operation. Aside from the ordinary raster scanning techniques, a common approach is to use light patterns codifying functions of a basis, such as Hadamard or Fourier components [26,27]. In this case, the image is obtained computationally by just an inverse transformation of the corresponding basis. The technique is very well adapted to apply compressive sensing algorithms by using different basis of functions, highly incoherent between them, to sample the object and to reconstruct the image [28].
Ghost imaging methods are similar to active single-pixel techniques, particularly in the case of computational ghost imaging. In these approaches the object is scanned with random light patterns and the image is obtained by correlation operations [29,30]. Active single-pixel imaging techniques are also related to dual photography, where the image is obtained by exploiting the Helmholtz reciprocity principle [31]. Temporal modulation imaging is based on similar principles. In this case, images are obtained by using a single-pixel detector and scanning the object with temporally multiplexing spatial frequency functions from orthogonal, time varying spatial line modulation gratings [32].
As we have mentioned above, the capabilities of single-pixel imaging systems have been verified in many different applications [39,40]. Also, it has been proposed to use LED arrays for single-pixel imaging in fluorescence application, achieving a very high pattern modulation frequency (250 kHz) although at a low resolution of 8 × 8 pixels [41]. A remarkable single-pixel imaging system to obtain 2D monochromatic images, working at a very high frequency with a resolution of 32 × 32 pixels, was developed very recently [42].
In this paper we propose a method to perform 2D color imaging with a single photodiode by using light structured illumination generated with a low-cost LED array. Moreover, the system is extended to 3D imaging by simply adding a low number of photodiodes in different locations. The LED array is used to generate a sequence of Hadamard patterns which are projected onto the object by a simple optical system. A field programmable gate array (FPGA) controls the LCD panel allowing us to obtain high refresh rates (up to 10 kHz). Color is obtained by generating a Hadamard pattern for each RGB chromatic component by controlling the color of the LEDs, which allows us to use a single monochromatic photodiode. Experimental results are provided for an LED frequency operation of 10 kHz.
Information of the 3D structure of the object is obtained by combining a low number of images generated from photodiodes located at different positions. We take advantage from the fact that shifting a photodiode in single-pixel imaging is equivalent to shifting the light source in conventional imaging [19,31]. We use a photometric stereo technique based on shape-from-shading [43], which obtains information about the 3D surface from a series of images of the object acquired while varying the illumination direction in the scene. There are two different methodologies to tackle the problem, called uncalibrated and calibrated photometric stereo. In this paper we use an uncalibrated method [44], which does not require to know the positions of the light sources in advance, allowing us to develop a very flexible 3D imaging system.
Single-pixel 3D imaging by using a color LED display
The layout of our single-pixel imaging system based on a color LED display is shown in Fig. 1. An LED panel controlled by an FPGA is used to generate a sequence of light patterns codifying the Hadamard functions. The FPGA also stores the binary patterns to be sequentially sent. The light patterns are projected onto the object by using a conventional optical system. The light reflected by the object is collected by a single photodiode and digitized by a DAQ system attached to a computer. To obtain a color image, a sequence of Hadamard patterns is projected for each RGB chromatic component. The image is generated in the computer by a simple Hadamard inverse transformation for each chromatic component.
Let us consider an LED panel containing a total of N light emitters. In mathematical terms, the image of the object provided by photodiode k is considered as a column vector, Ok, containing the N components of the object, in lexicographical order, which will be sampled by the light patterns. The light intensity measured by the photodiode can be written as: where H is the Hadamard matrix, with a size N × N, containing the N Walsh-Hadamard functions necessary to sample the object. Equation (1) is, in fact, the Hadamard transformation of the object onto the Hadamard frequency space. The object is then recovered by the inverse transformation 1 For the case of 3D imaging, a low number K of photodiodes is located at different positions around the object, as is shown in Fig. 1. Alternatively, a single photodiode is shifted to different positions and the sequence of patterns is sent again. As is well known, shifting the photodiode in single-pixel imaging techniques is equivalent to illuminate the object from a different direction in conventional imaging. This property has been used in dual photography [31] and in single-pixel 3D imaging by applying shape from shading algorithms [19]. In our system, with this approach we are able to obtain a set of different images to apply photometric stereo techniques without preliminary calibration. L ∈ . In our case, the LED projector is at a fixed position and the photodiodes are moved to K different positions to obtain the different images. If we assume that surfaces in the scene have a Lambertian behavior, the photodiode has a linear response, and assuming an orthographic projection, the object provided by our single-pixel camera in Fig. 1 from photodiode k, can be written as · In general, to obtain matrix G implies to find 9 parameters, for example by using information from pixels in which the relative value of the surface reflectance is constant or is known. The ambiguity can be further reduced to 3 parameters, the so-called generalized bas-relief ambiguity (GBR) [47], using an integrability constraint, which enforces that a consistent surface can be reconstructed from the estimated normal field. In the uncalibrated method used here, the GBR ambiguity is eliminated by exploiting points where the Lambertian reflection is maximal, called Lambertian diffuse reflectance maxima (LDRM) points [44]. These points are located approximately at the same pixel position of the image intensity maxima. It has been show that this method is tolerant to the presence of ambient illumination and shadows. Figure 2 shows a zenithal picture of the LED-based single-pixel camera configured by following the layout in Fig. 1. We use a 32 × 32 RGB LED matrix panel, with 4mm grid spacing (Adafruit). A simple lens, with a focal length f = 100 mm, located at a distance 35 cm far from the matrix, is used to project the patterns onto the object. The sensor is an amplified photodiode (Thorlabs PDA 100A-EC) and the FPGA module is an FPGA development KIT (Terasic, DE1-SoC board) with the Quartus Prime software used to codify the light patterns. The resolution of the Hadamard patterns is limited by that of the LED array. Therefore, the images reconstructed by our system have a resolution limited to 32 × 32 pixels. For the optical configuration in Fig. 2, the size of the scene illuminated by the light patterns is 4 × 4 cm. The pattern modulation frequency is limited by the maximum clock frequency of the LED array, which is 25MHz. Our 32 × 32 LED matrix is electronically composed by two matrices with 16 × 32 LEDs that can operate simultaneously. Then the maximum modulation frequency for patterns with 32 × 32 pixels is around 40 kHz if we take into account the multiplexing time and the control signals. However, the quality of the images in terms of the signal to noise ratio (SNR) was not satisfactory for frequencies larger than 10 kHz. Without applying CS, or similar optimization algorithms, we need 32 × 32 Hadamard patterns to obtain an image. Thus, at a pattern projection rate of 10 kHz, the image frame rate is almost 10 Hz.
Experimental results
We have performed two imaging experiments at a sampling rate of 10 kHz with the setup in Fig. 2. In the first one, we just record color images of the input scene. To this end, we send the full set of 1024 Walsh-Hadamard patterns associated to a space with 32 × 32 pixels, for each primary color, red, green and blue. The color image is reconstructed from the coefficients provided by the photodiode intensity signal by using Eq. (2). The result of this experiment is shown in Fig. 3. To quantify the quality of the image we measured the SNR for each chromatic channel RGB. The result of the averaged SNR for the three measurements is 53 dB. We measured also the SNR of the 2D monochromatic images used in the 3D imaging experiment shown in Fig. 4. As an example, the SNR of the central image in Fig. 4(a) is 62 dB. The SNR increases by using all LED colors at the same time to codify the light patterns. Fig. 3. 2D image of a color object obtained with the optical system shown in Fig. 2. The resolution is 32 × 32 pixels, the same as the number of elements in the LED array.
We can compare the resolution, modulation frequency, and cost of our single-pixel camera based on LEDs with those of light projection systems based on spatial light modulators. The resolution of our system is much lower than that provided by off-the-shelf video projectors based on digital micromirror device (DMD) or liquid crystal on silicon (LCOS) spatial light modulators, which is typically 1280 × 1024 pixels. However, the modulation frequency of these video projectors is limited to 60 Hz, which is two orders of magnitude lower than the frequency of the LED system with a similar cost. Some commercial video projectors allow to operate at higher frequencies by modulating binary patterns instead of three grey-level chromatic channels. In this case the modulation frequency can arrive to 1200 Hz but at a higher cost. Finally, scientific spatial light modulators based on DMDs, when controlled with high quality drivers, can arrive to frequencies as high as 22 kHz with high resolution and flexibility. However, with a cost higher in one order of magnitude. In our second experiment we obtain a 3D reconstruction of the surface of the object. As explained in the previous section, firstly we record a set of images of the object with our single-pixel camera in Fig. 2, by shifting the photodiode to different positions. In particular, we record 5 images by moving the photodiode in the same horizontal plane (see Fig. 1). This ensures us that we had enough LDRM points to obtain a robust GBR parameter estimation. The images are shown in Fig. 4(a). It is possible to see how the object appears illuminated from different directions. These images are used as inputs of the uncalibrated photometric stereo algorithm in [44]. The software code to apply this method is available in [48]. The resulting 3D shape is shown in Fig. 4(b) and 4(c). The axis labels in Fig. 4(b) show the x, y, and z coordinates of the object surface in mm. Please note that the surrounding of the figure looks uniform blue because the shape-from-shading algorithm do not provide information for low values of the intensity. Figure 4(c) shows just a 3D representation of the surface.
To estimate the quality of our result, we have measured relative depths at different locations in the object with a Vernier caliper appreciating 0.05 mm. The object has a size of 32.5 × 37.2 × 27.1 mm. The relative error of our surface, averaged from measurements in three locations of the object, is 3.1 ± 0.2%. Also, we have estimated the error of an uncalibrated shape-from-shading method in the bibliography performed with a conventional camera. In the experiment by Hayakawa [46], the average error is estimated to be 2.2%, which is of the same order of magnitude of our result. Finally, we have compared the result of our single-pixel uncalibrated 3D imaging method with the single-pixel calibrated one described in [19]. The root mean square error (RMSE) of our measurements compared to those performed with the caliper is 0,17 mm for an object with a total depth of 27.1 mm, while the RMSE of the measurements performed with the system in [19] compared with those provided by a stereo photogrammetric camera system is 4 mm for an object with a total depth of 250 mm. The RMSE error is lower in our case, but it should be taken into account that we measured the error only from a few measurements, while in [19] the comparison was performed for all the surface. Besides, in our case the total depth is much lower. The relative error in terms of the total depth is of the same order of magnitude in both cases.
Conclusions
We have proposed a method to perform 2D color imaging and 3D imaging with a single photodiode by using light structured illumination generated with an LED array. Our system allows us to record low-resolution images of color objects with a resolution limited to 32 × 32 pixels, the number of LEDs of the array. Despite the maximum modulation frequency of the LED panel is 40 kHz, we need to operate at a frequency of 10 kHz to obtain a satisfactory SNR. Therefore, the image frame rate is almost 10 Hz. With this frequency we obtained 2D images with SNR equal to 53 dB for color images and 62 dB for monochromatic ones.
The imaging system has a cost similar to off-the shelf video projectors based on DMD or LCOS spatial light modulators, which have a higher resolution but are limited to operate at lower frequencies in single-pixel imaging applications. Spatial light modulators based on DMDs and controlled with high quality drivers provide both larger resolutions and modulation frequencies but at a higher cost. The resolution of our system is low because of the limited size and pitch of the LED array. It could be possible to improve the resolution by combining several LED panels operating in parallel at the same frequency. In our case, however, the characteristics of our FPGA board would limit the total number of panels to three.
For 3D imaging we have used a non-calibrated photometric stereo technique avoiding a precise calibration of the position of the photodiodes. The accuracy of the depth profile was determined by direct comparison with measurements using a Vernier caliper. We estimated a depth RMSE of 0,17 mm and an averaged relative error of 3.1%.
We have used the Walsh-Hadamard basis for the sampling operation, but other different basis such as Fourier or Morlet wavelet could also be used. In this prove of concept experiment we did not performed CS but we plan to do it in the next applications of the camera. By using CS and optimizing the system to obtain higher frequency rates, we expect to get color 2D images at a frame rate of 40 fps. | 5,878.6 | 2018-06-06T00:00:00.000 | [
"Engineering",
"Physics"
] |
Chemically Reduced Graphene Oxide-Reinforced Poly(Lactic Acid)/Poly(Ethylene Glycol) Nanocomposites: Preparation, Characterization, and Applications in Electromagnetic Interference Shielding
In this study, a nanocomposite of reduced graphene oxide (RGO) nanofiller-reinforcement poly(lactic acid) (PLA)/poly(ethylene glycol) (PEG) matrix was prepared via the melt blending method. The flexibility of PLA was improved by blending the polymer with a PEG plasticizer as a second polymer. To enhance the electromagnetic interference shielding properties of the nanocomposite, different RGO wt % were combined with the PLA/PEG blend. Using Fourier-transform infrared (FT-IR) spectroscopy, field emission scanning electron microscopy (FE-SEM) and X-ray diffraction, the structural, microstructure, and morphological properties of the polymer and the RGO/PLA/PEG nanocomposites were examined. These studies showed that the RGO addition did not considerably affect the crystallinity of the resulting nanomaterials. Thermal analysis (TGA) reveals that the addition of RGO highly improved the thermal stability of PLA/PEG nanocomposites. The dielectric properties and electromagnetic interference shielding effectiveness of the synthesized nanocomposites were calculated and showed a higher SE total value than the target value (20 dB). On the other hand, the results showed an increased power loss by increasing the frequency and conversely decreased with an increased percentage of filler.
Introduction
Electromagnetic interference (EMI) is one of the most undesirable byproducts of the rapid growth of telecommunication devices and high-frequency electronic systems. Any device that utilizes, distributes, processes, or transmits any form of electrical energy is likely to interfere with nearby equipment or systems' operation and emit electromagnetic signals [1]. Negative effects on the health of human being might also be a result of such phenomenon. There has been a great deal of effort invested in the reduction of electromagnetic pollution using EMI shielding materials [2], where signals are attenuated by the shielding materials through absorption and/or reflection of the radiation power [3]. Magnetic materials and metallic structures have been used traditionally for EMI the most feasible method. In this method, chemically reduced graphene oxide (RGO) (graphene), was prepared through oxidizing graphite in the presence of oxidants and strong acids, followed by graphene oxide (GO) reduction using one of the reduction chemical methods [31]. On its surface, the residual oxygen functional groups were retained by the resulting RGO, which will greatly influence its properties and that of its nanocomposites [32].
The electromagnetic properties and EMI shielding effectiveness of nanocomposites depend on the reflection from the material's surface, absorption of the EM energy, and propagation paths of the EM wave, which are determined by the nature, shape, size, and microstructure of the fillers [33]. There are some studies that have investigated and contributed to the knowledge of the effect of the particle size of the graphene oxide on the electromagnetic properties, where the properties can be tuned by controlling size of the graphene particles. This phenomenon may be interpreted as follows: the interfacial polarization is dominant in the heterogeneous structures, which contributes to the various electromagnetic properties [34]. Also, defects of graphene oxide act as dipolar polarization centers, leading to various dielectric relaxations in different frequency ranges. With decreasing nanoparticle size, larger interface area and more defects exist in the nanoparticles, which enhance the electromagnetic properties of the nanoparticles [35].
In this work, the RGO/PLA/PEG nanocomposites' preparation was carried out by using the technique called melt blending using a Brabender internal mixer. The properties of the RGO/PLA/PEG nanocomposites, including the crystallization behaviour, the functional groups of the nanocomposites, the structure of the composite, and the thermal behaviour of the composites, are thoroughly investigated. Furthermore, the effectiveness of the EMI shielding properties of nanocomposites over the X band (8)(9)(10)(11)(12) GHz) is also investigated. It was found that the applied frequency and the filler concentrations have affected all the essential properties of the RGO/PLA/PEG nanocomposites.
Materials
The main material utilized in this study is the PLA polymer matrix pellets with a density of 1.24 g/cm 3 prepared through oxidizing graphite in the presence of oxidants and strong acids, followed by graphene oxide (GO) reduction using one of the reduction chemical methods [31]. On its surface, the residual oxygen functional groups were retained by the resulting RGO, which will greatly influence its properties and that of its nanocomposites [32]. The electromagnetic properties and EMI shielding effectiveness of nanocomposites depend on the reflection from the material's surface, absorption of the EM energy, and propagation paths of the EM wave, which are determined by the nature, shape, size, and microstructure of the fillers [33]. There are some studies that have investigated and contributed to the knowledge of the effect of the particle size of the graphene oxide on the electromagnetic properties, where the properties can be tuned by controlling size of the graphene particles. This phenomenon may be interpreted as follows: the interfacial polarization is dominant in the heterogeneous structures, which contributes to the various electromagnetic properties [34]. Also, defects of graphene oxide act as dipolar polarization centers, leading to various dielectric relaxations in different frequency ranges. With decreasing nanoparticle size, larger interface area and more defects exist in the nanoparticles, which enhance the electromagnetic properties of the nanoparticles [35].
In this work, the RGO/PLA/PEG nanocomposites' preparation was carried out by using the technique called melt blending using a Brabender internal mixer. The properties of the RGO/PLA/PEG nanocomposites, including the crystallization behaviour, the functional groups of the nanocomposites, the structure of the composite, and the thermal behaviour of the composites, are thoroughly investigated. Furthermore, the effectiveness of the EMI shielding properties of nanocomposites over the X band (8-12 GHz) is also investigated. It was found that the applied frequency and the filler concentrations have affected all the essential properties of the RGO/PLA/PEG nanocomposites.
Materials
The main material utilized in this study is the PLA polymer matrix pellets with a density of 1.24 g/cm 3 (Grade 4060D) from Nature Work LLC (Minnetonka, MN, USA). Low molecular weight polyethylene glycol (PEG) (Mn = 200 g/mol) was acquired from Sigma-Aldrich (St. Louis, MO, USA). Reduced Graphite Oxide powder was synthesized via a chemical method in the lab. NH3 was supplied by Sigma Aldrich (Sarasota, FL, USA). The chemical structures of the materials under test are illustrated in Figure 1.
Preparation of Reduced Graphene Oxide (RGO)
Reduced graphene oxide (RGO) powder was synthesized by a chemical method with an average size of 60 nm. The chemical method of manufacturing and preparing the RGO powder involved two major steps. The first step is the synthesis of Graphite Oxide (GO) using the Staudenmaier Method [35]. The second step performs the reduction process of GO to RGO. About 400 mg of the obtained GO was placed in a cellulose extraction thimble (30 by 100 mm) and was then placed in the Soxhlet extraction unit. Approximately 150 mL of 30% Ammonia solution (NH 3 ) was used as a reducing agent. The heating temperature was set at 90 • C, and the investigated exposure period was 5 h when the GO powder had direct contact with the ammonia vapor as well as the condensed liquid. The procedure of manufacturing and preparation is shown in Figure 2.
Preparation of Reduced Graphene Oxide (RGO)
Reduced graphene oxide (RGO) powder was synthesized by a chemical method with an average size of 60 nm. The chemical method of manufacturing and preparing the RGO powder involved two major steps. The first step is the synthesis of Graphite Oxide (GO) using the Staudenmaier Method [35]. The second step performs the reduction process of GO to RGO. About 400 mg of the obtained GO was placed in a cellulose extraction thimble (30 by 100 mm) and was then placed in the Soxhlet extraction unit. Approximately 150 mL of 30% Ammonia solution (NH3) was used as a reducing agent. The heating temperature was set at 90 °C, and the investigated exposure period was 5 h when the GO powder had direct contact with the ammonia vapor as well as the condensed liquid. The procedure of manufacturing and preparation is shown in Figure 2.
Preparation of RGO/PLA/PEG Nanocomposites
To remove the moisture content of PLA during blending, both PLA pellets and RGO powder were dried at 80 °C for 12 h in a vacuum oven before processing [36]. The RGO/PLA/PEG nanocomposites' preparation was carried out by using the technique called melt blending for 10 min using Brabender internal mixer (GmbH & Co. KG, Duisburg, Germany) at 170 °C with 60 rpm rotor speed [37]. The PLA to PEG weight ratio was kept constant at 90/10 wt/wt. The RGO was added as a filler at a different percentage with the PLA/PEG blend as presented in Table 1. The dispersion of RGO nanofiller within the PLA/PEG polymer leads to network formation as a result of more physical contacts between the particles of filler and matrix inside the RGO/PLA/PEG nanocomposites. The structure of RGO/PLA/PEG nanocomposites is schematically illustrated in Figure 3.
Preparation of RGO/PLA/PEG Nanocomposites
To remove the moisture content of PLA during blending, both PLA pellets and RGO powder were dried at 80 • C for 12 h in a vacuum oven before processing [36]. The RGO/PLA/PEG nanocomposites' preparation was carried out by using the technique called melt blending for 10 min using Brabender internal mixer (GmbH & Co. KG, Duisburg, Germany) at 170 • C with 60 rpm rotor speed [37]. The PLA to PEG weight ratio was kept constant at 90/10 wt/wt. The RGO was added as a filler at a different percentage with the PLA/PEG blend as presented in Table 1. The dispersion of RGO nanofiller within the PLA/PEG polymer leads to network formation as a result of more physical contacts between the particles of filler and matrix inside the RGO/PLA/PEG nanocomposites. The structure of RGO/PLA/PEG nanocomposites is schematically illustrated in Figure 3.
Preparation of Reduced Graphene Oxide (RGO)
Reduced graphene oxide (RGO) powder was synthesized by a chemical method with an average size of 60 nm. The chemical method of manufacturing and preparing the RGO powder involved two major steps. The first step is the synthesis of Graphite Oxide (GO) using the Staudenmaier Method [35]. The second step performs the reduction process of GO to RGO. About 400 mg of the obtained GO was placed in a cellulose extraction thimble (30 by 100 mm) and was then placed in the Soxhlet extraction unit. Approximately 150 mL of 30% Ammonia solution (NH3) was used as a reducing agent. The heating temperature was set at 90 °C, and the investigated exposure period was 5 h when the GO powder had direct contact with the ammonia vapor as well as the condensed liquid. The procedure of manufacturing and preparation is shown in Figure 2.
Preparation of RGO/PLA/PEG Nanocomposites
To remove the moisture content of PLA during blending, both PLA pellets and RGO powder were dried at 80 °C for 12 h in a vacuum oven before processing [36]. The RGO/PLA/PEG nanocomposites' preparation was carried out by using the technique called melt blending for 10 min using Brabender internal mixer (GmbH & Co. KG, Duisburg, Germany) at 170 °C with 60 rpm rotor speed [37]. The PLA to PEG weight ratio was kept constant at 90/10 wt/wt. The RGO was added as a filler at a different percentage with the PLA/PEG blend as presented in Table 1. The dispersion of RGO nanofiller within the PLA/PEG polymer leads to network formation as a result of more physical contacts between the particles of filler and matrix inside the RGO/PLA/PEG nanocomposites. The structure of RGO/PLA/PEG nanocomposites is schematically illustrated in Figure 3. The molding of the obtained composite into sheets of thickness 1 mm was carried out by hot pressing for 10 min at 170 • C with a pressure of 110 k/bar, followed by room temperature cooling. The prepared plates were used for further characterization. As for the electromagnetic properties of nanocomposites, the dough was molded into rectangular-shaped specimens with a thickness of 6 mm by hot pressing. The preparation temperature was 170 • C and the pressure force was kept at 110 k/bar for 10 min, after which the specimens were allowed to cool down for 10 min at a force of 110 k/bar to a temperature of 50 • C.
The rectangular mold specimen was placed between the two rectangular waveguides to minimize the air gap between waveguide walls and the border of the specimen. Figure 4a,b illustrates the process used in fabricating the substrates composites and the specimens where the mixture was poured into rectangular aluminum molds of 6 mm thickness, in preparation for the electromagnetic measurements of composites. The molding of the obtained composite into sheets of thickness 1 mm was carried out by hot pressing for 10 min at 170 °C with a pressure of 110 k/bar, followed by room temperature cooling. The prepared plates were used for further characterization. As for the electromagnetic properties of nanocomposites, the dough was molded into rectangular-shaped specimens with a thickness of 6 mm by hot pressing. The preparation temperature was 170 °C and the pressure force was kept at 110 k/bar for 10 min, after which the specimens were allowed to cool down for 10 min at a force of 110 k/bar to a temperature of 50 °C.
The rectangular mold specimen was placed between the two rectangular waveguides to minimize the air gap between waveguide walls and the border of the specimen. Figure 4a, b illustrates the process used in fabricating the substrates composites and the specimens where the mixture was poured into rectangular aluminum molds of 6 mm thickness, in preparation for the electromagnetic measurements of composites.
Morphological Characterization
FE-SEM images provide in-depth information which reveals the characteristic 3-D manifestation necessary for understanding the morphology of the cross-sectional surface of a sample. Accordingly, the surface morphology of RGO/PLA/PEG nanocomposites and the dispersion of RGO nanoparticles in the PLA/PEG matrix were studied using FE-SEM (FEI Quanta 200 SEM, Yuseong, Daejeon, Korea) at a fixed voltage of 10 kV. The specimens were dried for 45 min before being coated with gold particles using a SEM coating unit (Baltic SC030 sputter coater, Yuseong, Daejeon, Korea).
X-Ray Diffraction
The measurement of X-ray diffraction was performed by using an X-ray diffractometer (XRD, XD-3, Cu Ka radiation) under ambient conditions with a Lynx Eye detector using a Bruker diffractometer ((Yuseong, Daejeon, Korea) over 2θ range of 5°-80°. The used X-ray beam was Cu-Kα
Morphological Characterization
FE-SEM images provide in-depth information which reveals the characteristic 3-D manifestation necessary for understanding the morphology of the cross-sectional surface of a sample. Accordingly, the surface morphology of RGO/PLA/PEG nanocomposites and the dispersion of RGO nanoparticles in the PLA/PEG matrix were studied using FE-SEM (FEI Quanta 200 SEM, Yuseong, Daejeon, Korea) at a fixed voltage of 10 kV. The specimens were dried for 45 min before being coated with gold particles using a SEM coating unit (Baltic SC030 sputter coater, Yuseong, Daejeon, Korea).
X-Ray Diffraction
The measurement of X-ray diffraction was performed by using an X-ray diffractometer (XRD, XD-3, Cu Ka radiation) under ambient conditions with a Lynx Eye detector using a Bruker diffractometer ((Yuseong, Daejeon, Korea) over 2θ range of 5 • -80 • . The used X-ray beam was Cu-Kα where the average crystallite size is given as D (in nm), the shape factor (normally 0.9 for cubic) is k, the full-width half maximum intensity measured diffraction line broadening is β (FWHM data converted to radians), Bragg's diffraction angle is θ, and the X-ray's wavelength is λ.
Thermal Stability (TGA and DTG) Properties
Thermogravimetric analysis (TGA) analysis is a technique in which a measurement of the mass of a substance is carried out depending on a temperature or time function while the material undergoes a controlled temperature program. TGA is also useful for compositional analysis of multicomponent materials and used to examine the kinetics of the physiochemical processes occurring in the sample, whereas the Derivative thermogravimetric (DTG) thermograms were utilized for composites' weight-loss study. The investigation of the RGO/PLA nanocomposites' thermal stability was performed using a TGA (1600LF, Shanghai Mettler Toledo Co. Ltd., Shanghai, China). The RGO/PLA specimens, weighing 5-10 mg each, were heated from 50 • C to 600 • C at a heating rate of 10 • C/min −1 under a nitrogen atmosphere with a flow rate of 20 mL/min. The instrument was computer controlled while Pyris software was used for calculations. The instrument was calibrated using the Curie temperatures of five different metal standards.
Fourier Transform Infrared (FT-IR) Analysis
A Perkin-Elmer FT-IR 1650 spectrophotometer (Waltham, MA, USA) was used for the FT-IR characterization, identification of chemical bonds in a molecule, and determination of the functional groups in the net and RGO/PLA/PEG nanocomposites at different percentage of RGO. The FT-IR spectroscopy involves the collection of information on IR absorption and analyzing them in spectrum form by correlating the frequencies of IR radiation absorption ("peaks" or "signals") directly to the bonds present in the compound under investigation. Using a KBr disk method, the FT-IR spectra tests were carried out in the 400 to 4000 cm −1 wavenumber range, at room temperature.
Electrical Properties
The dielectric properties and Scattering (S-parameters) reflection (S11) and transmission (S21) coefficients were measured using a transmission line technique (rectangular waveguide) by commercial measurement software on the Agilent N5230A PNA-L network analyzer system which includes the Agilent 85071E, 85701B software package, respectively [38] (Agilent Tech, 2010, Keysight Technologies, Santa Rosa, CA, USA), in the 8-12 GHz frequency range. The rectangular waveguide connected to the vector network analyzer (VNA) has been calibrated by full two-port, Thru-Reflect-Line (TRL) [39].
EMI Shielding Effectiveness (EMI SE)
The electromagnetic interference shielding effectiveness (EMI SE) is defined as the ability of a shield material to reduce the electromagnetic field. It can also be defined by the ratio of incoming to outgoing power. The total EMI SE of a material is specified to take place by the three mechanisms namely: absorption shielding effectiveness (SE A ), reflection shielding effectiveness (SE R ), and multiple internal reflections (SE M ). A schematic mechanism of EMI SE is depicted in Figure 5. The scattering parameters (S 11 , S 12 , S 21 , and S 22 ) were obtained as shielding components for the conductive nanocomposites. The shielding effectiveness for the nanocomposites has been calculated by using Equation (2) [40]. The absorption mechanism is associated with the dielectric/magnetic polarization or energy dissipation [41]. The reflection, which is considered as the primary mechanism of shielding, involves the interactions between the electromagnetic fields and the charge carriers, such as electrons and holes. It can also be related with the impedance mismatch between the air and an absorber. Internal multiple reflections are related to the reflection between the opposite faces of a material. Therefore, in order to barricade the electromagnetic wave, the shielding material should either absorb or reflect the wave. In general, EMI SE is expressed in decibel (dB) units.
The EMI SE of a material can be expressed using the following equation [40]: where P in and P tr are power of incident and transmitted EM waves, respectively. In the case of a thick shield, the SE M can be neglected due to high absorption loss. When, for the second time, the wave reaches the second boundary, its amplitude becomes negligible, due to the fact that it has passed through the shield's thickness three times by then. The EMI SE equation can be written as [42]: Therefore, experimental absorption and reflection losses can be expressed as [43]: where the S 11 and S 21 scattering parameters are coefficients of reflection and transmission, respectively. Equation (2) [40]. The absorption mechanism is associated with the dielectric/magnetic polarization or energy dissipation [41]. The reflection, which is considered as the primary mechanism of shielding, involves the interactions between the electromagnetic fields and the charge carriers, such as electrons and holes. It can also be related with the impedance mismatch between the air and an absorber.
Internal multiple reflections are related to the reflection between the opposite faces of a material. Therefore, in order to barricade the electromagnetic wave, the shielding material should either absorb or reflect the wave. In general, EMI SE is expressed in decibel (dB) units.
The EMI SE of a material can be expressed using the following equation [40]: Where Pin and Ptr are power of incident and transmitted EM waves, respectively. In the case of a thick shield, the SEM can be neglected due to high absorption loss. When, for the second time, the wave reaches the second boundary, its amplitude becomes negligible, due to the fact that it has passed through the shield's thickness three times by then. The EMI SE equation can be written as [42]: Therefore, experimental absorption and reflection losses can be expressed as [43]: where the S11 and S21 scattering parameters are coefficients of reflection and transmission, respectively.
Field Emission-Scanning Electron Microscopy Results
The results of FE-SEM were used for the study of the fractured tensile specimens' surface morphology and to visualize qualitatively the RGOs' dispersion state in the matrix of the PLA/PEG mixture. The surface micrographs of RGO/PLA/PEG nanocomposites, PLA/PEG mixture, and the neat PLA are shown in Figure 6a-f. Figure 6a shows a typical tensile fracture surface of PLA.
Field Emission-Scanning Electron Microscopy Results
The results of FE-SEM were used for the study of the fractured tensile specimens' surface morphology and to visualize qualitatively the RGOs' dispersion state in the matrix of the PLA/PEG Figure 6a shows a typical tensile fracture surface of PLA. Thermoplastic organizing of the PLA after melting is accredited to real melting processes due to its brittle behavior at room temperature, as shown in the Figure [35]. Illustrated in Figure 6b, the PLA/PEG mixture is subjected to the process of deformation and a few long threads of deformed material are discernible on the fracture surface of the sample. Figure 6c clearly shows the layered, porous, wrinkled silk-like morphology; the layered structure was formed by RGO particles and continually cross-linked in a flaky textured form, as described by [44]. Thermoplastic organizing of the PLA after melting is accredited to real melting processes due to its brittle behavior at room temperature, as shown in the Figure [35]. Illustrated in Figure 6b, the PLA/PEG mixture is subjected to the process of deformation and a few long threads of deformed material are discernible on the fracture surface of the sample. Figure 6c clearly shows the layered, porous, wrinkled silk-like morphology; the layered structure was formed by RGO particles and continually cross-linked in a flaky textured form, as described by [44]. To understand the dispersion of RGO in the PLA/PEG mixture, the wrinkled morphology of RGO is assumed to be necessary in the mechanical interlocking with the polymer matrix, which helped to build a strong interfacial interactions and subsequently efficient stress transfer across the interface [45]. The RGO particles are located on the surface of PLA/PEG matrix and trapped between the matrix like a sandwich in some areas, indicating an electrostatic attraction between RGO powder and matrix, which may contribute to the microwave absorption [46].
Good composite uniformity indicates a good dispersion degree of RGO and thus results in good thermal stability and tensile properties. On the other hand, after increasing the percentage of RGO (2.4% and 4%) in the PLA/PEG mixture, the dispersion of RGO powder in the matrix in Figure 6e,f showed obvious differences. Figure 6e showed that the RGO powder was dispersed in the PLA/PEG mixture as merged particles of a large size. Figure 6f showed that the RGO powder was dispersed in the PLA/PEG mixture in the form of agglomerates. The outlines of the PLA/PEG mixture and RGO nanoparticles are clearly observable. In addition, it can be seen that the small RGO particles have a propensity to aggregate due to their inherent properties. However, it is considered that suitable RGO nanoparticle content could cause uniform distribution at the matrix and suppress the aggregation. To understand the dispersion of RGO in the PLA/PEG mixture, the wrinkled morphology of RGO is assumed to be necessary in the mechanical interlocking with the polymer matrix, which helped to build a strong interfacial interactions and subsequently efficient stress transfer across the interface [45]. The RGO particles are located on the surface of PLA/PEG matrix and trapped between the matrix like a sandwich in some areas, indicating an electrostatic attraction between RGO powder and matrix, which may contribute to the microwave absorption [46].
Fourier Transform Infrared (FT-IR) Analysis
Good composite uniformity indicates a good dispersion degree of RGO and thus results in good thermal stability and tensile properties. On the other hand, after increasing the percentage of RGO (2.4% and 4%) in the PLA/PEG mixture, the dispersion of RGO powder in the matrix in Figure 6e,f showed obvious differences. Figure 6e showed that the RGO powder was dispersed in the PLA/PEG mixture as merged particles of a large size. Figure 6f showed that the RGO powder was dispersed in the PLA/PEG mixture in the form of agglomerates. The outlines of the PLA/PEG mixture and RGO nanoparticles are clearly observable. In addition, it can be seen that the small RGO particles have a propensity to aggregate due to their inherent properties. However, it is considered that suitable RGO nanoparticle content could cause uniform distribution at the matrix and suppress the aggregation. Figure 6 presents the FT-IR spectra for pure material, PLA/PEG mixture, and RGO/PLA/PEG nanocomposites. The PLA spectrum shows four main regions corresponding to (1100-1000) cm −1 , (1500-1400) cm −1 , (1750-1745) cm −1 , and (3000-2850) cm −1 , respectively. The PEG spectrum shows a clear peak at the absorption bands corresponding to 3441 cm −1 , 2878 cm −1 , 1464, 1343, and 1279 cm −1 , respectively [47]. Figure 7 shows the RGO spectrum, with four main regions of absorption bands of the RGO powder, which correspond to 1039.46 cm −1 , 1388.03 cm −1 , 1494.24 cm −1 , 1590.97 cm −1 , 3223.46 cm −1 , and 3389.86 cm −1 , respectively [35]. The distinguishing peaks responsible for -C=O stretching, -CH stretching, -C-O stretching, as well as C-H bending were clearly observed for all RGO/PLA/PEG nanocomposites over the spectra and no new peaks were formed. Therefore, a conclusion can be made that no chemical interaction took place in the polymer matrices with RGO addition. This is expected since RGO does not have any functional groups with which to form a strong interface with a polymer matrix. Thus, any change in the nanocomposites' properties is the effect of the physical interaction only between the PLA/PEG mixture and the RGO powder. It was also clear from the results that no considerable change was observed in the peak positions of the RGO/PLA/PGE nanocomposites compared to pristine PLA/PGE mixture.
Fourier Transform Infrared (FT-IR) Analysis
Polymers 2019, 11, x FOR PEER REVIEW 9 of 21 Figure 6 presents the FT-IR spectra for pure material, PLA/PEG mixture, and RGO/PLA/PEG nanocomposites. The PLA spectrum shows four main regions corresponding to (1100-1000) cm −1 , (1500-1400) cm −1 , (1750-1745) cm −1 , and (3000-2850) cm −1 , respectively. The PEG spectrum shows a clear peak at the absorption bands corresponding to 3441 cm −1 , 2878 cm −1 , 1464, 1343, and 1279 cm −1 , respectively [47]. Figure 7 shows the RGO spectrum, with four main regions of absorption bands of the RGO powder, which correspond to 1039.46 cm −1 , 1388.03 cm −1 , 1494.24 cm −1 , 1590. 97 cm −1 , 3223.46 cm −1 , and 3389.86 cm −1 , respectively [35]. The distinguishing peaks responsible for -C=O stretching, -CH stretching, -C-O stretching, as well as C-H bending were clearly observed for all RGO/PLA/PEG nanocomposites over the spectra and no new peaks were formed. Therefore, a conclusion can be made that no chemical interaction took place in the polymer matrices with RGO addition. This is expected since RGO does not have any functional groups with which to form a strong interface with a polymer matrix. Thus, any change in the nanocomposites' properties is the effect of the physical interaction only between the PLA/PEG mixture and the RGO powder. It was also clear from the results that no considerable change was observed in the peak positions of the RGO/PLA/PGE nanocomposites compared to pristine PLA/PGE mixture. Figure 8 illustrates the X-ray diffraction pattern of the PLA/PEG mixture, RGO powder, and RGO/PLA/PEG nanocomposites at different percentages of RGO filler. For the PLA/PEG mixture, a broad amorphous peak from PLA was observed which is the typical peak for any given amorphous structure around 17.48°, this is in agreement with [48]. Figure 7 confirms the X-ray diffraction pattern of the RGO powder, showing a good crystallinity. The curve indicates a series of diffraction peaks at 2θ = 17.85°, 38.57°, 42.23°, 44.87°, and 73.02°, which correspond to the (d002), (d100), (d101), (d102), and (d004) planes, respectively [35]. Normally, the disappearance of the characteristic peak for RGO powder in the RGO/PLA/PEG nanocomposites can be correlated to fully exfoliated RGO powder in the polymer matrix. This indicates that the accumulated layers of RGO have exceeded high shearing during melt-mixing. The RGO/PLA/PEG nanocomposites at all percentages showed almost the same absorption peaks as pristine PLA. This means that there is no new bond created or strong chemical interaction occurring between the matrix and the nanocomposites. Figure 8 illustrates the X-ray diffraction pattern of the PLA/PEG mixture, RGO powder, and RGO/PLA/PEG nanocomposites at different percentages of RGO filler. For the PLA/PEG mixture, a broad amorphous peak from PLA was observed which is the typical peak for any given amorphous structure around 17.48 • , this is in agreement with [48]. Figure 7 confirms the X-ray diffraction pattern of the RGO powder, showing a good crystallinity. The curve indicates a series of diffraction peaks at 2θ = 17.85 • , 38.57 • , 42.23 • , 44.87 • , and 73.02 • , which correspond to the (d002), (d100), (d101), (d102), and (d004) planes, respectively [35]. Normally, the disappearance of the characteristic peak for RGO powder in the RGO/PLA/PEG nanocomposites can be correlated to fully exfoliated RGO powder in the polymer matrix. This indicates that the accumulated layers of RGO have exceeded high shearing during melt-mixing. The RGO/PLA/PEG nanocomposites at all percentages showed almost the same absorption peaks as pristine PLA. This means that there is no new bond created or strong chemical interaction occurring between the matrix and the nanocomposites.
Thermogravimetric Analysis (TGA)
During its service life, an EMI shielding material may be subject to high-temperature conditions. In this study, the thermal stability of the composites was analyzed. A feasible EMI shield should have suitable thermal stability, so that it may perform well at elevated temperatures [49]. Reduced graphene oxide has been applied widely as a filler to improve the thermal stability of the polymer matrix. In this research, thermogravimetric analysis (TGA) and the curves of the derivative (DTG) thermogravimetric are used to gain information about the extent and nature of the material degradation under test and RGO/PLA/PEG nanocomposites. The detailed experimental results are shown in Figure 9 and Table 2. The TGA results illustrated the PEG/PLA matrix showed a lower thermal stability compared to pure PLA. The reduction in the thermal stability of PLA is mainly a result of PEG's presence as a plasticizer. PEG causes a reduction in thermal stability by its action to intersperse itself around polymers and by breaking the interaction between the PLA and PEG polymers, which is predicted by the gel theory of plasticization and lubricity theory [50].
Thermogravimetric Analysis (TGA)
During its service life, an EMI shielding material may be subject to high-temperature conditions. In this study, the thermal stability of the composites was analyzed. A feasible EMI shield should have suitable thermal stability, so that it may perform well at elevated temperatures [49]. Reduced graphene oxide has been applied widely as a filler to improve the thermal stability of the polymer matrix. In this research, thermogravimetric analysis (TGA) and the curves of the derivative (DTG) thermogravimetric are used to gain information about the extent and nature of the material degradation under test and RGO/PLA/PEG nanocomposites. The detailed experimental results are shown in Figure 9 and Table 2. The TGA results illustrated the PEG/PLA matrix showed a lower thermal stability compared to pure PLA. The reduction in the thermal stability of PLA is mainly a result of PEG's presence as a plasticizer. PEG causes a reduction in thermal stability by its action to intersperse itself around polymers and by breaking the interaction between the PLA and PEG polymers, which is predicted by the gel theory of plasticization and lubricity theory [50].
Thermogravimetric Analysis (TGA)
During its service life, an EMI shielding material may be subject to high-temperature conditions. In this study, the thermal stability of the composites was analyzed. A feasible EMI shield should have suitable thermal stability, so that it may perform well at elevated temperatures [49]. Reduced graphene oxide has been applied widely as a filler to improve the thermal stability of the polymer matrix. In this research, thermogravimetric analysis (TGA) and the curves of the derivative (DTG) thermogravimetric are used to gain information about the extent and nature of the material degradation under test and RGO/PLA/PEG nanocomposites. The detailed experimental results are shown in Figure 9 and Table 2. The TGA results illustrated the PEG/PLA matrix showed a lower thermal stability compared to pure PLA. The reduction in the thermal stability of PLA is mainly a result of PEG's presence as a plasticizer. PEG causes a reduction in thermal stability by its action to intersperse itself around polymers and by breaking the interaction between the PLA and PEG polymers, which is predicted by the gel theory of plasticization and lubricity theory [50]. The thermal degradation temperature of pure RGO powder, pure PLA, and RGO/PLA/PEG nanocomposites is between 227 °C and 400 °C. The thermal decomposition curves of RGO was at 227 The thermal degradation temperature of pure RGO powder, pure PLA, and RGO/PLA/PEG nanocomposites is between 227 • C and 400 • C. The thermal decomposition curves of RGO was at 227 • C, while PLA was at 356 • C. The thermal decomposition curves of RGO/PLA/PEG nanocomposites at different percentage of filler shifted towards higher temperatures, compared to the pure PLA/PEG matrix with increasing RGO content. The 50% weight loss temperatures (T 50% • C) of pure PLA was 365 • C, while at 0.8%, 2.4%, and 4% RGO content, they were 330 • C, 340 • C, and 350 • C, respectively. Figure 9 and Table 2 show that the T d-max , which represents maximum decomposition temperature, also shifted towards higher temperatures with increasing RGO content, and the Td-max = 386.167 • C of the composites was increased with 4.0 wt% RGO filler compared to pure PLA polymer. This suggests that RGO nanoparticles can control the thermal stability of the RGO/PLA/PEG nanocomposites as well as T 50% and T d-max . At the beginning, the degradation temperature of the composites can be ascribed to the early decomposition of RGO in the matrix. At 50% weight loss, the thermal decomposition temperature of the composites was found to have improved. Generally, RGO incorporation into the PLA/PEG mixture enhanced the thermal stability by acting as a superior insulator and mass transport barrier to the volatile products generated during decomposition. Many researchers have also demonstrated that the incorporation of graphene or its chemical derivatives could enhance the thermal stability of PLA at extremely low loading contents [51]. It has to be noted that the PLA/PEG blends and RGO/PLA/PEG nanocomposites showed lower degradation temperatures compared to the neat PLA due to the thermal decomposition of the polymer matrix [52]. Table 2 shows weight loss (%) at T d-max temperatures. It was observed that the weight loss (%) of pure PLA, PLA/PEG, RGO, and its nanocomposite (0.8% RGO, 2.4% RGO, and 4% RGO) at 400 • C were 97%, 95%, 27%, 92%, 90%, and 87% respectively. The delayed degradation of the PLA/PEG chain with increasing concentration of RGO led to the enhanced thermal stability of the nanocomposites. It was observed that the weight loss % of RGO at T d-max was 27% because the thermal decomposition of RGO starts slightly below 400 • C, and maximum decomposition occurs above 600 • C, leading to the enhancement of thermal stability. Figure 10 shows the derivative thermograms (DTG) of pure PLA, PLA/PEG mixture, and RGO/PLA/PEG nanocomposites at different percentages of filler, while the RGO performs as a heat barrier, which enhances the overall thermal stability of the nanocomposites. It also assists in the formation of ash after thermal decomposition. The RGO shifts the decomposition to a higher temperature at early stages of thermal decomposition. The improvement in thermal stability can be attributed to the "tortuous path" effect of RGO, which retards the escape of volatile degradation products. The presence of high RGO loading or well-dispersed RGO in the polymer matrix will force the degradation products to go through the more tortuous path and hence enhanced the thermal stability. Similar results have been reported for other graphene-based nanocomposites [53].
The Dielectric Properties of the Composites
The permittivity (ε) of a material determines the material's response to the electric field component of the electromagnetic wave. Permittivity is determined by the complex term ε = ε′ − jε″. Insulating polymers have low permittivity due to the small degree of polarization of the macromolecules. Addition of conductive fillers to polymers can lead to considerable improvement to the low permittivity of the matrix [54] since the polarization of filler and polarization of filler/polymer interface (interfacial polarization) can contribute significantly to the overall polarization of the composite. When a current flow across the interface of two materials, it can cause accumulation of charges at the interface due to the difference in the materials' relaxation times and consequently increases its permittivity [55]. Figure 11a, b shows the real and imaginary part of permittivity of the RGO/PLA/PEG nanocomposites versus RGO loading. It is evident that the permittivity is very sensitive to RGO loading. Both the real and imaginary parts of the permittivity are found to have increased with increasing RGO concentration. The real part increases reasonably from 2.75 to 3.79 as the RGO loading is raised from 0 to 4 wt% and the imaginary part increases from 0.093 to 0.50 for the same increment in filler content. This significant improvement in both dielectric constant and loss factor is the result of an increase in conductivity and dipole moment of RGO/PLA/PEG nanocomposites due to addition of conductive nanoparticles [56].
The Dielectric Properties of the Composites
The permittivity (ε) of a material determines the material's response to the electric field component of the electromagnetic wave. Permittivity is determined by the complex term ε = ε − jε". Insulating polymers have low permittivity due to the small degree of polarization of the macromolecules. Addition of conductive fillers to polymers can lead to considerable improvement to the low permittivity of the matrix [54] since the polarization of filler and polarization of filler/polymer interface (interfacial polarization) can contribute significantly to the overall polarization of the composite. When a current flow across the interface of two materials, it can cause accumulation of charges at the interface due to the difference in the materials' relaxation times and consequently increases its permittivity [55]. Figure 11a,b shows the real and imaginary part of permittivity of the RGO/PLA/PEG nanocomposites versus RGO loading. It is evident that the permittivity is very sensitive to RGO loading. Both the real and imaginary parts of the permittivity are found to have increased with increasing RGO concentration. The real part increases reasonably from 2.75 to 3.79 as the RGO loading is raised from 0 to 4 wt% and the imaginary part increases from 0.093 to 0.50 for the same increment in filler content. This significant improvement in both dielectric constant and loss factor is the result of an increase in conductivity and dipole moment of RGO/PLA/PEG nanocomposites due to addition of conductive nanoparticles [56].
The permittivity of the composites is related to the absorption of an electromagnetic wave along with the dielectric constant and the thickness of materials. It is anticipated that the conductive network formed by the RGO interacts with the entering power signal and assists the movement of an electron within the composites, thus being responsible for the absorption of incident power. Also, the interaction between the RGO and PLA/PEG matrix contributes to the movement of electrons in the composites. It is expected that the conducting RGO renders multiple interfaces that increases the reflection and greatly barricades the electromagnetic wave inside the composites. The fine dispersed RGO particles facilitate easy movement of free electrons inside the insulating polymer matrix, even at low filler loading. Figure 11. The frequency dependence of (a) dielectric constant (ε′), (b) loss factor (ε″) at various RGO loadings.
The permittivity of the composites is related to the absorption of an electromagnetic wave along with the dielectric constant and the thickness of materials. It is anticipated that the conductive network formed by the RGO interacts with the entering power signal and assists the movement of an electron within the composites, thus being responsible for the absorption of incident power. Also, the interaction between the RGO and PLA/PEG matrix contributes to the movement of electrons in the composites. It is expected that the conducting RGO renders multiple interfaces that increases the reflection and greatly barricades the electromagnetic wave inside the composites. The fine dispersed RGO particles facilitate easy movement of free electrons inside the insulating polymer matrix, even at low filler loading.
EMI Shielding Effectiveness
The EMI SE of RGO/PLA/PEG nanocomposites was determined in the frequency range of [8][9][10][11][12] GHz which is illustrated in Figure 12a-d. The EMI SE values of the RGO/PLA/PEG nanocomposites were calculated using Equations (2) -(5), respectively. The PLA/PEG matrix is transparent to the electromagnetic radiation and does not show any EMI shielding efficiency due to the very low permittivity of pure PLA. The SEtotal of the nanocomposites increases upon increasing RGO loading due to the enhancement of permittivity (ε′ and ε″). With the increase in filler loading, the EMI shielding efficiency of RGO/PLA/PEG nanocomposites increases. The EMI SE value of 22.5 dB was Figure 11. The frequency dependence of (a) dielectric constant (ε ), (b) loss factor (ε") at various RGO loadings.
EMI Shielding Effectiveness
The EMI SE of RGO/PLA/PEG nanocomposites was determined in the frequency range of 8-12 GHz which is illustrated in Figure 12a-d. The EMI SE values of the RGO/PLA/PEG nanocomposites were calculated using Equations (2)-(5), respectively. The PLA/PEG matrix is transparent to the electromagnetic radiation and does not show any EMI shielding efficiency due to the very low permittivity of pure PLA. The SE total of the nanocomposites increases upon increasing RGO loading due to the enhancement of permittivity (ε and ε"). With the increase in filler loading, the EMI shielding efficiency of RGO/PLA/PEG nanocomposites increases. The EMI SE value of 22.5 dB was recorded at 0.8 wt % RGO loading. The high value of EMI SE obtained at low loading of functionalized RGO can be attributed to the fine distribution and dispersion of RGO in the PLA/PEG matrix, thereby forming an interconnected network. The electrical behaviors of the material are crucial for EMI shielding efficiency since it is responsible for interacting with the electromagnetic wave. The total EMI SE is influenced by the mesh size and the amount of mobile charge carriers provided by the filler network in the composites. The EMI SE results of these nanocomposites showed higher values than the minimum value of EMI SE of the shielding required for practical applications, which are usually rated around 20 dB. matrix, thereby forming an interconnected network. The electrical behaviors of the material are crucial for EMI shielding efficiency since it is responsible for interacting with the electromagnetic wave. The total EMI SE is influenced by the mesh size and the amount of mobile charge carriers provided by the filler network in the composites. The EMI SE results of these nanocomposites showed higher values than the minimum value of EMI SE of the shielding required for practical applications, which are usually rated around 20 dB. Figure 12a shows that the reflection loss (SER) increased with frequency, from 5.56 to 18.96 dB at low percentages of filler (0.8% RGO). Also, SER increases gradually from 4.52 to 11.65 dB at high percentage of filler (4% RGO). The absorption loss SEA decreased from 6.30 to 3.61 dB at the low percentage of filler (0.8% RGO), and the (SEA) decreased gradually from 6.62 to 3.50 dB at a high percentage of filler (4% RGO) as shown in Figure 12b. Figure 12c indicates the proportional relation between the SE total values to the frequency. The inverse proportional effect of RGO loading on SER and SEtotal values is shown in Figure 12d, where the RGO increment reflect a reduction of SER and SEtotal. While the SEA values are increasing to RGO loading increment. The mean EMI SE values of the composites with 0.8%, 1.60%, 2.4%, 3.2%, and 4% mass fractions of RGO are presented in Table 3. Figure 12a shows that the reflection loss (SE R ) increased with frequency, from 5.56 to 18.96 dB at low percentages of filler (0.8% RGO). Also, SE R increases gradually from 4.52 to 11.65 dB at high percentage of filler (4% RGO). The absorption loss SE A decreased from 6.30 to 3.61 dB at the low percentage of filler (0.8% RGO), and the (SE A ) decreased gradually from 6.62 to 3.50 dB at a high percentage of filler (4% RGO) as shown in Figure 12b. Figure 12c indicates the proportional relation between the SE total values to the frequency. The inverse proportional effect of RGO loading on SE R and SE total values is shown in Figure 12d, where the RGO increment reflect a reduction of SE R and SE total. While the SE A values are increasing to RGO loading increment. The mean EMI SE values of the composites with 0.8%, 1.60%, 2.4%, 3.2%, and 4% mass fractions of RGO are presented in Table 3. The EMI SE results of the RGO/PLA/PEG nanocomposites in the current work were compared to the previously reported composites prepared with different mixing approaches, different sample thickness, and similar or even lower/higher conductive filler loading, as listed in Table 4. The results of current RGO/PLA/PEG nanocomposites exhibit an efficient EMI SE even when compared to the best values for different polymer/conductive filler nanocomposites. The minimum value of EMI SE of the shielding required for practical application is usually considered to be~20 dB.
Power Loss and Affective Absorbance A eff
While total shielding effectiveness is an important parameter commonly used to quantify the efficiency of a shielding material, it does not provide information on the contributions of each of the shielding mechanisms. To determine the influence of RGO powder loading on the reflection and absorption of the nanocomposites, power balance calculations were performed at various frequencies, which can be expressed by following the generalized equation; The power loss results of the RGO/PLA/PEG nanocomposites displayed in Figure 13 show that the low values of power loss in the case of RGO/PLA/PEG nanocomposites containing high amounts of RGO are due to the low power transmitted into the sample as a result of a very good reflection [64] at the sample's surface. A better understanding of RGO/PLA/PEG nanocomposites' potential to absorb electromagnetic radiation can be obtained by evaluating their effective absorbance. The intensity of an EM wave inside the material after primary reflection is based on quantity (1-R), which can be used for adjustment of absorbance (A) to yield effective absorbance from the formula [65]: A eff determines what percentage of the power entering the material has been absorbed. Therefore, it is essential to differentiate between the absolute amount of power absorbed by a shielding material and its potential for absorption. Materials that have high reflection can be used for EMI shielding in cases where reflection is of no concern. On the other hand, materials that have high absorption potential (SE A ) but also reflect a significant fraction of the incident power at their surface may be used as radiation absorbers if their reflection could be reduced via impedance matching at the surface as in structure-engineered shielding, such as multilayer structures or via foaming, etc. [66]. Figure 14 shows the variation of effective absorbance values to the frequency in the RGO/PLA/PEG nanocomposites at different RGO loading. Where the A eff values of all the samples decrease to the frequency rise, whereas, the higher A eff values have been observed at the higher RGO loading in the PLA/PEG polymer matrix.
at the sample's surface. A better understanding of RGO/PLA/PEG nanocomposites' potential to absorb electromagnetic radiation can be obtained by evaluating their effective absorbance. The intensity of an EM wave inside the material after primary reflection is based on quantity (1-R), which can be used for adjustment of absorbance (A) to yield effective absorbance from the formula [65]: Aeff (%) = (1 -R − T)/(1 − R) × 100 (7) Aeff determines what percentage of the power entering the material has been absorbed. Therefore, it is essential to differentiate between the absolute amount of power absorbed by a shielding material and its potential for absorption. Materials that have high reflection can be used for EMI shielding in cases where reflection is of no concern. On the other hand, materials that have high absorption potential (SEA) but also reflect a significant fraction of the incident power at their surface may be used as radiation absorbers if their reflection could be reduced via impedance matching at the surface as in structure-engineered shielding, such as multilayer structures or via foaming, etc. [66]. Figure 14 shows the variation of effective absorbance values to the frequency in the RGO/PLA/PEG nanocomposites at different RGO loading. Where the Aeff values of all the samples decrease to the frequency rise, whereas, the higher Aeff values have been observed at the higher RGO loading in the PLA/PEG polymer matrix.
Conclusions
Reduced graphene oxide (RGO)-reinforced polylactic acid (PLA)/polyethylene glycol (PEG) blended nanocomposites were prepared by the melt blending method. FE-SEM micrographs showed good dispersion of RGO nanoparticles in the PLA/PEG matrix at low concentrations, while at higher loadings, RGO was found to become physically in contact, forming a conductive passage inside the matrix. X-ray diffraction shows that the addition of RGO did not considerably affect the crystallinity of the resulting nanocomposite materials. The improved thermomechanical properties of the composites yielded to a good adhesion between the RGO nanoparticles and the PLA/PEG matrix. The RGO acts as a heat barrier, which enhances the overall thermal stability of the polymer nanocomposites and assists in the formation of ash after thermal decomposition. The dielectric properties of the PLA/PEG matrix were enhanced remarkably with the addition of RGO. The EMI shielding properties of the synthesized composites were tested and the composites showed that the SE total value is higher than the target value of 20 dB. Reflection was found to be the dominant shielding mechanism for RGO/PLA/PEG nanocomposites over the X band. However, the contribution of absorption to SEtotal increased as the RGO loading was decreased. The effective
Conclusions
Reduced graphene oxide (RGO)-reinforced polylactic acid (PLA)/polyethylene glycol (PEG) blended nanocomposites were prepared by the melt blending method. FE-SEM micrographs showed good dispersion of RGO nanoparticles in the PLA/PEG matrix at low concentrations, while at higher loadings, RGO was found to become physically in contact, forming a conductive passage inside the matrix. X-ray diffraction shows that the addition of RGO did not considerably affect the crystallinity of the resulting nanocomposite materials. The improved thermomechanical properties of the composites yielded to a good adhesion between the RGO nanoparticles and the PLA/PEG matrix. The RGO acts as a heat barrier, which enhances the overall thermal stability of the polymer nanocomposites and assists in the formation of ash after thermal decomposition. The dielectric properties of the PLA/PEG matrix were enhanced remarkably with the addition of RGO. The EMI shielding properties of the synthesized composites were tested and the composites showed that the SE total value is higher than the target value of 20 dB. Reflection was found to be the dominant shielding mechanism for RGO/PLA/PEG nanocomposites over the X band. However, the contribution of absorption to SE total increased as the RGO loading was decreased. The effective absorbance of the matrix increased with loading increase. The results showed increased power loss with an increase in the frequency and conversely decreased with an increased percentage of filler. Effective absorbance of the materials increased with loading increase. The materials showed high absorption potential (SE A ), but also reflected a significant fraction of the incident power at their surface. Thus, the materials may be used as radiation absorbers if their reflection could be reduced via impedance matching at the surface, as in structure-engineered shielding, such as multilayer structures or via foaming. | 12,293.4 | 2019-04-01T00:00:00.000 | [
"Materials Science",
"Physics",
"Engineering"
] |
Intra-Patient Heterogeneity of Circulating Tumor Cells and Circulating Tumor DNA in Blood of Melanoma Patients
Despite remarkable progress in melanoma therapy, the exceptional heterogeneity of the disease has prevented the development of reliable companion biomarkers for the prediction or monitoring of therapy responses. Here, we show that difficulties in detecting blood-based markers, like circulating tumor cells (CTC), might arise from the translation of the mutational heterogeneity of melanoma cells towards their surface marker expression. We provide a unique method, which enables the molecular characterization of clinically relevant CTC subsets, as well as circulating tumor DNA (ctDNA), from a single blood sample. The study demonstrates the benefit of a combined analysis of ctDNA and CTC counts in melanoma patients, revealing that CTC subsets and ctDNA provide synergistic real-time information on the mutational status, RNA and protein expression of melanoma cells in individual patients, in relation to clinical outcome.
Introduction
Malignant Melanoma is the deadliest of all skin cancers and accounted for more than 59,000 deaths worldwide in 2015 [1]. In recent years, systemic treatment of metastatic melanoma has been transformed. Improved understanding of the genetic landscape of melanoma led to the development of BRAF and MEK inhibitors for patients with BRAF mutated tumors [2,3]. However, the frequently profound response to BRAF/MEK inhibition is transient in about 50% of all cases. Additional therapeutic options were derived from insights into the molecular controls of the immune system. CTLA-4 and PD-1/PDL-1 neutralizing antibodies are used independent of mutational status, leading to a durable response. However, this only occurs in a subset of patients [4,5].
Hence, reliable biomarkers, which allow the prediction of therapeutic response and/or the development of therapeutic resistance as early as possible, are urgently needed. To date, tissue biopsies have predominantly been utilized to achieve this goal. However, repeated biopsies to study the frequently adapting heterogeneous tumor cell populations in melanoma are invasive, difficult to obtain and may not represent the entire molecular tumor profile [6][7][8].
Circulating tumor cells (CTCs), as well as circulating tumor DNA (ctDNA), are shed into the bloodstream from either primary or metastatic lesions. Serial analysis of liquid biopsies might provide a dynamic and minimally invasive option to screen the pathological characteristics of the entire tumor, based on a simple blood withdrawal [9,10]. PCR-based studies on melanoma-associated antigens (MAAs, e.g., MART-1, MAGE-A3, PAX3, and GM2/GD2)-that are not present in leukocytes-have shown that their presence was correlated with an advanced patient stage, as well as decreased disease-free and overall survival rates, in several studies [11][12][13][14]. In addition, an increased quantity of ctDNA was found to be prognostic in melanoma patients and might provide useful information on the mutational status of the disease [15,16]. Furthermore, ctDNA might serve as a surrogate marker of tumor burden in metastatic melanoma patients [17]. Although promising, the information presented by these assays is still limited in its capability to advise therapeutic decisions, since it is based on the analysis of a pooled cell fraction (healthy and tumor tissue).
However, the enrichment of melanoma CTCs has been challenging, mainly due to their large molecular heterogeneity (e.g., surface marker expression and cellular size). To date, the use of either marker or size-based enrichment methods leads to the loss of surface marker negative or small CTCs, respectively. Interestingly, evidence is accumulating that the mutational heterogeneity of melanoma cells might influence the cellular expression of surface markers as well as their cellular volume. For example, activation of the RAS/RAF pathway drives the expression of CSPG4, the most commonly used surface protein, to enrich melanoma CTCs [18,19]. BRAF-inhibition decreases the cellular volume of enlarged BRAF-mutated melanoma cells in a glucose dependent manner [20]. The identification of all CTC subpopulations might, therefore, be pivotal for the correct stratification of patients and subsequent therapeutic decisions. In addition to diagnostic applications, the detailed analysis of CTC subpopulations may yield new insights into the process of melanoma metastasis.
Here, we show that the cell surface marker expression of melanoma cells depends on their mutational status. We provide a novel enrichment approach, which allows the isolation and complete molecular analysis of different CTC subpopulations and ctDNA analysis from one blood sample. In addition, we demonstrate how combined CTC and ctDNA analyses can reveal synergistic information, which is potentially relevant for personalized therapy in metastatic melanoma.
RAS/RAF Activating Mutations Lead to A Distinct Melanoma Marker Expression Pattern
Activation of the RAS/RAF pathway has been suggested to increase the expression of the melanoma surface marker CSPG4 in neural cells [18]. Here, we tested whether this is also the case for the expression of the previously described [21,22] melanoma-specific genes, by analyzing The Cancer Genome Atlas (TCGA) (https://www.cancergenome.nih.gov/) database. Melanomas containing RAS/RAF activating mutations were compared to melanomas without RAS/RAF activating mutations.
In total, 16 melanoma marker genes (S100A1, ABCB5, CDH19, MIA, SLC26A2, MCAM, S100A2, S100P, MAGEA4, TFAP2C, SFRP1, SERPINA3, CSPG4, TYRP1, IL13RA2, S100A7A) were highly expressed in melanomas with activating mutations ( Figure 1A,C). We next analyzed whether those differentially expressed genes could be detected on single cell level within each group. For this purpose, data from 2056 single melanoma cells, derived from the metastatic tumors of 19 melanoma patients, were analyzed [21]. t-SNE clustering clearly demonstrated that cells harboring a RAS/RAF activating Cancers 2019, 11, 1685 3 of 17 mutation can be differentiated from cells without activating mutations, based on melanoma marker genes ( Figure 1B). In total, 13 melanoma marker genes were specifically increased in the RAS/RAF cohort in both the bulk tumor (TCGA) and single melanoma cell cohort [21] ( Figure 1C). These markers include the two most commonly used melanoma surface markers for CTC enrichment, CSPG4 and MCAM ( Figure 1D).
Cancers 2019, 11, x 3 of 17 melanoma marker genes ( Figure 1B). In total, 13 melanoma marker genes were specifically increased in the RAS/RAF cohort in both the bulk tumor (TCGA) and single melanoma cell cohort [21] ( Figure 1C). These markers include the two most commonly used melanoma surface markers for CTC enrichment, CSPG4 and MCAM ( Figure 1D).
A Combined Enrichment Approach Allows the Detection of CTC Subpopulations
We reasoned that differential marker expression between mutational melanoma subsets might distort clinical decisions based on CTC counts and characteristics, depending on the utilized enrichment method. To identify a method which would (i) allow the detection of both marker positive and marker negative cells and (ii) enable a thorough molecular analysis of the isolated single cells (e.g., DNA, RNA, immunocytochemistry), we tested two marker dependent, three marker independent and one combined approach, for CTC detection (Figure 2A). The recovery rate for each (A) Volcano plot of differentially regulated genes in the TCGA data set, comparing melanomas containing BRAF/NRAS mutations (RAS/RAF activating) and not BRAF/NRAS mutated melanomas (other). Significantly reduced genes (<lgFC − 1) are depicted in violet; significantly increased genes (>lgFC 1) in yellow (B) t-SNE plot, based on melanoma markers expressed on single cells (Tirosh et al.), derived from BRAF/NRAS mutated tumors and not BRAF/NRAS mutated tumors. (C) Venn diagram of all differentially expressed genes between RAS/RAF activating and other tumors of the TCGA and Tirosh cohorts and specific melanoma marker genes. (D) Significantly differentially expressed melanoma marker genes in the Tirosh data set. * = p < 0.05.
A Combined Enrichment Approach Allows the Detection of CTC Subpopulations
We reasoned that differential marker expression between mutational melanoma subsets might distort clinical decisions based on CTC counts and characteristics, depending on the utilized enrichment method. To identify a method which would (i) allow the detection of both marker positive and marker negative cells and (ii) enable a thorough molecular analysis of the isolated single cells (e.g., DNA, RNA, immunocytochemistry), we tested two marker dependent, three marker independent and one combined approach, for CTC detection (Figure 2A). The recovery rate for each method was determined by spiking 25 individually micro-manipulated cells of an RAF mutated (SKMEL28) and non-RAS/RAF mutated (MeWo) cell each into 7.5 mL blood from healthy donors. The recovery rate varied between 36% (Cellsearch ® ) and 82% (combined approach) ( Figure 2B).
An additional obstacle for the translation of each method into clinical practice is provided by the number of contaminating leukocytes, represented by the number of histological slides needed to analyze cellular output after each enrichment method (1 million resulting cells were mounted onto each slide for the subsequent stain and detection of the cells). For complete analysis of a Leucosep ® enriched sample, an average of 12 slides was needed, whereas only one slide was necessary after CellSearch ® enrichment ( Figure 2B). As expected, subsequent staining of the cells for CSPG4/MCAM and S100 (as a marker for surface marker negative cells) revealed that marker independent methods tend to isolate a more representative cellular population than marker dependent methods ( Figure 2C).
For the following analysis, we decided to use the combined approach, which allowed a good resolution of the melanoma cell subpopulations, and yielded a high recovery rate, a low amount of contaminating leukocytes, and a complete molecular analysis of CTCs.
Targeted Sequencing Reveals Mutational CTC Subclones
We collected whole blood samples from 84 melanoma patients receiving current standards of clinical care, to determine if CTC subtypes can be used to support clinical diagnostics. Patients presenting with stage I-IV cutaneous, acral, amelanotic, lentigo, desmoplastic or uveal melanoma were included. Patients were between 21-88 years old and received treatment, including An additional obstacle for the translation of each method into clinical practice is provided by the number of contaminating leukocytes, represented by the number of histological slides needed to analyze cellular output after each enrichment method (1 million resulting cells were mounted onto each slide for the subsequent stain and detection of the cells). For complete analysis of a Leucosep ® enriched sample, an average of 12 slides was needed, whereas only one slide was necessary after CellSearch ® enrichment ( Figure 2B). As expected, subsequent staining of the cells for CSPG4/MCAM and S100 (as a marker for surface marker negative cells) revealed that marker independent methods tend to isolate a more representative cellular population than marker dependent methods ( Figure 2C).
For the following analysis, we decided to use the combined approach, which allowed a good resolution of the melanoma cell subpopulations, and yielded a high recovery rate, a low amount of contaminating leukocytes, and a complete molecular analysis of CTCs.
Targeted Sequencing Reveals Mutational CTC Subclones
We collected whole blood samples from 84 melanoma patients receiving current standards of clinical care, to determine if CTC subtypes can be used to support clinical diagnostics. Patients presenting with stage I-IV cutaneous, acral, amelanotic, lentigo, desmoplastic or uveal melanoma were included. Patients were between 21-88 years old and received treatment, including chemotherapy, targeted therapies and immunotherapy. Some patients have been followed up for more than three years ( Figure 3A).
Overall, 32% (27 patients, Supplementary Table S1) of all patients were CTC positive. An increase in CTC-positive patients was detected with increased tumor staging ( Figure 3H). The mean number of CTCs was 4.85 and the median was 3.0.
Patients with stage I or II disease harbored CTCs which were either enriched by positive selection (stage I and II) or showed CTCs in both positive and size dependent enrichment (stage II), hinting at a high expression of cell surface markers. Stage IV patients showed positivity in all enrichment approaches, which was also reflected in treatment naïve patients ( Figure 3I). In accordance with the finding that RAS/RAF activating mutations result in a higher expression of melanoma surface markers ( Figure 1), we detected more RAS/RAF mutated cells in the cellular subpopulation, which was enriched by positive selection, in comparison to the Parsortix™ enriched subpopulation.
Patient 1 presented with an NRAS Q61K primary melanoma and a positive sentinel lymph node biopsy at the time of the first liquid biopsy ( Figure 3B-G). At the same time, a CSPG4/MCAM positive CTC was detected, containing a NRAS Q61PL mutation. At week 4, the patient developed a lymph node metastasis and satellite metastasis, that were treated by surgical resection. After 62 weeks, we detected two marker positive CTCs, which contained a NRAS Q61RL and a BRAF V600E mutation and a surface marker negative CTC, without any RAS/RAF driver mutation. At week 74, the patient clinically relapsed (subcutaneous metastasis, SC). Between weeks 78-82, the SC metastasis was treated by radiotherapy, resulting in a complete response. At week 98, CTC analysis revealed one marker negative and RAF/RAS negative CTC. The patient relapsed at week 128, and was treated successfully with Pembrolizumab (PD-1 inhibitor). Note that CT scans, performed at weeks 12, 50, and 79, did not show any sign of metastasis or progression and LDH (lactate dehydrogenase) levels did not reflect the recurrence of the metastasis. S100 levels were elevated throughout the complete follow up period, thus limiting its predictive power, although this could possibly be helpful in the detection of minimal residual disease. chemotherapy, targeted therapies and immunotherapy. Some patients have been followed up for more than three years ( Figure 3A). Overall, 32% (27 patients, Supplementary Table S1) of all patients were CTC positive. An increase in CTC-positive patients was detected with increased tumor staging ( Figure 3H). The mean number of CTCs was 4.85 and the median was 3.0.
Patients with stage I or II disease harbored CTCs which were either enriched by positive selection (stage I and II) or showed CTCs in both positive and size dependent enrichment (stage II), hinting at a high expression of cell surface markers. Interestingly, surface marker expression seems to be reduced or lost in higher stage patients, as reflected in the exclusive detection of CTCs by the size dependent Parsortix TM approach. Stage IV patients showed positivity in all enrichment approaches, which was also reflected in treatment naïve patients ( Figure 3I). In accordance with the finding that RAS/RAF activating mutations result in a higher expression of melanoma surface markers ( Figure 1), we detected more RAS/RAF mutated cells in the cellular subpopulation, which was enriched by positive selection, in comparison to the Parsortix TM enriched subpopulation.
Patient 1 presented with an NRAS Q61K primary melanoma and a positive sentinel lymph node biopsy at the time of the first liquid biopsy ( Figure 3B-G). At the same time, a CSPG4/MCAM positive CTC was detected, containing a NRAS Q61PL mutation. At week 4, the patient developed a lymph node metastasis and satellite metastasis, that were treated by surgical resection. After 62 weeks, we detected two marker positive CTCs, which contained a NRAS Q61RL and a BRAF V600E mutation and a surface marker negative CTC, without any RAS/RAF driver mutation. At week 74, the patient clinically relapsed (subcutaneous metastasis, SC). Between weeks 78-82, the SC metastasis was treated by radiotherapy, resulting in a complete response. At week 98, CTC analysis revealed one marker negative and RAF/RAS negative CTC. The patient relapsed at week 128, and was treated successfully with Pembrolizumab (PD-1 inhibitor). Note that CT scans, performed at weeks 12, 50, and 79, did not show any sign of metastasis or progression and LDH (lactate dehydrogenase) levels did not reflect the recurrence of the metastasis. S100 levels were elevated throughout the complete follow up period, thus limiting its predictive power, although this could possibly be helpful in the detection of minimal residual disease.
RNA Expression Pattern on Selected CTCs
Since we aimed to develop a CTC enrichment method which allows for the thorough analysis of melanoma CTC subpopulations, we tested the feasibility of the method for the analysis of RNA expression on single cell level. We first performed a pathway analysis in the TCGA cohort, comparing RAS/RAF mutated versus non-mutated melanomas. Using Cytoscape and ClueGO, a Go-term analysis of significantly upregulated genes (lgFC >2, adj. p-value 0.05) in the RAS/RAF cohort was performed. As expected, Go-terms, which are associated with the MAPK pathway, were significantly enriched (included in the "positive regulation of cell communication" cluster) (Supplementary Figure S1A). Interestingly, cell chemotaxis, which plays a major role in metastasis, was overrepresented as well. We next enriched CTCs from Patient 4 (BRAF V600E , stage IV) within 4 hours of blood withdrawal, using the combined approach, but refraining from the use of fixation methods. In total, three marker
RNA Expression Pattern on Selected CTCs
Since we aimed to develop a CTC enrichment method which allows for the thorough analysis of melanoma CTC subpopulations, we tested the feasibility of the method for the analysis of RNA expression on single cell level. We first performed a pathway analysis in the TCGA cohort, comparing RAS/RAF mutated versus non-mutated melanomas. Using Cytoscape and ClueGO, a Go-term analysis of significantly upregulated genes (lgFC > 2, adj. p-value < 0.05) in the RAS/RAF cohort was performed. As expected, Go-terms, which are associated with the MAPK pathway, were significantly enriched (included in the "positive regulation of cell communication" cluster) (Supplementary Figure S1A). Interestingly, cell chemotaxis, which plays a major role in metastasis, was overrepresented as well. We next enriched CTCs from Patient 4 (BRAF V600E , stage IV) within 4 h of blood withdrawal, using the combined approach, but refraining from the use of fixation methods. In total, three marker positive and two marker negative cells were detected. An analysis of genes present in the Go-term "regulation of cell motility" by qRT-PCR, showed that marker positive cells indeed showed the enrichment of these genes (Supplementary Figure S1B).
CTCs and ctDNA Provide Synergistic Clinical Information
Information regarding how useful ctDNA might be for the stratification of melanoma patients, and whether ctDNA provides additional or congruent information in comparison to CTCs, is still sparse, and was therefore scrutinized. The detection of ctDNA against the normally occurring background of cell-free DNA is challenging. One possible solution might be the characterization of ctDNA fragment size. ctDNA has been reported to be overrepresented in the fraction below 150 bp [23]. Since it is conceivable that the amount of recovered ctDNA tumor depends on the tumor burden, we compared the number of CTCs, ctDNA > 150 bp and ctDNA < 150 bp in patients, with regards to the Breslow thickness of the primary tumor, or the existence of a lymph node or systemic metastasis ( Figure 4A-C).
CTC counts did not dramatically change between tumors below 1 mm and between 1-2 mm, however, an increase in the average number of detected cells was seen in tumors with a Breslow thickness above 2.1 mm, which was further increased in patients with systemic disease. Total ctDNA showed only a slight increase with increasing tumor thickness. ctDNA concentration below 150 bp allowed a comparable discrimination between primary tumors below and above 2 mm. Thus, both CTC count and ctDNA < 150 bp appear to be a promising tool to predict tumor burden in our cohort.
Patient 2 was a metastatic melanoma patient, without any detected driver mutations in the primary tumor ( Figure 4D). At the time of first analysis, the patient presented a metastasis in the bone and the suprarenal gland (SG). At the same time, three CTCs were detected, and a ctDNA < 150 bp concentration of 0.9 ng/mL was measured. Both CTCs and ctDNA did not contain any driver mutations. The patient received Pembrolizumab, starting in week 4, and showed a partial response (bone). Meanwhile, the patient developed a pancreatic lesion. The PET-CT (positron emission tomography-computed tomography) confirmed a pancreatitis (possible side effect of Pembrolizumab), without any sign of new metastasis. At week 16, zero CTCs were detected. Two weeks later, we were able to detect five CTCs in this now-untreated metastatic melanoma patient, whereas the concentration of ctDNA < 150 bp was reduced, compared to the initial values. The PET-CT from the same day showed a progress of the SG metastasis and a new bone metastasis. Targeted sequencing of four out of the five CTCs revealed two BRAF V600E and one EGFR l491M mutation. At week 28, no CTCs were detected; however, ctDNA < 150 bp increased to~2 ng/mL and confirmed the development of a BRAF V600E mutated tumor. The patient progressed in week 36 (SG metastasis). Note that, in this case, LDH and S100 levels were poor markers for disease progression. S100, however, was dramatically increased at week 36.
Patient 3 was diagnosed with a BRAF V600K positive melanoma in 2012. After a recurrence in 2013 and 2014, and LN metastasis and lung metastasis in 2015, treatment with Dabrafenib (BRAF inhibitor) and Trametinib (MEK inhibitor) resulted in a complete remission. One year later, in weeks 0, 5 and 31, one CTC was detected at each time point ( Figure 4E). Note that, at week 11, the patient was diagnosed with schwannoma. Targeted sequencing showed BRAF mutations in all detected CTCs. A shift from BRAF V600K to BRAF V600E and later BRAF V600K plus MAP2K1 P124S was found. At week 38, no CTCs were detected. In comparison, ctDNA < 150 bp was elevated at the time point of the initial blood draw, decreased at 5 weeks, and increased again at week 31. No mutations were found at either 0 or 5 weeks. At week 31, a DPH3 mutation was detected. The patient relapsed at week 118, with an upper arm metastasis.
Melanomas have been known to quickly adapt their mutational pattern, in response to environmental and therapeutic pressure. Here, we tested whether mutations found in liquid biopsies of metastatic patients differed from the mutational status of the tissue derived from the primary tumor (reported by the department of pathology). The initial mutational status was recovered in 47.61-70.58% of all cases. Importantly, novel driver mutations were detected in 29.42-52.39% of all samples ( Figure 4F).
CTCs, ctDNA<150 bp and LDH Predict Clinical Outcome
We next analyzed whether stratification of patients by the existence of CTCs or a ctDNA concentration higher than 2 ng/mL for ctDNA >150 bp, or 0.5 ng/mL for ctDNA <150bp, would predict patient survival. The cut-off values were chosen based on values detected in blood from healthy donors. Kaplan-Meier curves demonstrate that patients with detectable CTC (≥1 CTC), ctDNA
CTCs, ctDNA < 150 bp and LDH Predict Clinical Outcome
We next analyzed whether stratification of patients by the existence of CTCs or a ctDNA concentration higher than 2 ng/mL for ctDNA > 150 bp, or 0.5 ng/mL for ctDNA < 150 bp, would predict patient survival. The cut-off values were chosen based on values detected in blood from healthy donors.
Kaplan-Meier curves demonstrate that patients with detectable CTC (≥1 CTC), ctDNA < 150 bp and LDH positive patients show a worse outcome than marker negative patients ( Figure 5A-E). In a cox-proportional hazards regression analysis, adjusted for stage, age, gender and treatment, the hazard ratio for LDH was calculated to be 5.07, followed by ctDNA < 150 bp (4.21) and CTCs (3.96). ctDNA > 150 bp and S100 were not found to significantly alter the HR in melanoma patients ( Figure 5F). Note that Patients 2 and 3 ( Figure 4D-E) were selected based on interesting and representative clinical courses, and are not representative of the HR values calculated in Figure 5F. <150bp and LDH positive patients show a worse outcome than marker negative patients ( Figure 5A-E). In a cox-proportional hazards regression analysis, adjusted for stage, age, gender and treatment, the hazard ratio for LDH was calculated to be 5.07, followed by ctDNA <150 bp (4.21) and CTCs (3.96). ctDNA >150 bp and S100 were not found to significantly alter the HR in melanoma patients ( Figure 5F). Note that Patients 2 and 3 ( Figure 4D-E) were selected based on interesting and representative clinical courses, and are not representative of the HR values calculated in Figure 5F.
Discussion
Our work showed that surface marker expression on melanoma cells is dependent on their mutational status. Furthermore, we demonstrated that a combined analysis of ctDNA and CTCs predicted relapse earlier than imaging, and was more accurate than serum LDH or S100 in a subset of patients. Interestingly, we were able to detect "private" mutations on CTCs and ctDNA, that were not revealed in the random bulk analysis of the primary tumor.
In the present study, we have analyzed melanoma-associated cell surface markers in relation to the mutational status of the melanoma cells. We found a larger proportion of surface marker-positive cells (e.g., CSPG4/MCAM) in the RAS/RAF-mutated cohort compared to the non-RAF/RAS mutated cohort. Thus, we conclude that the commonly employed enrichment of CTCs based on surface marker expression might be biased and could lead to the loss of subsets of tumor cells lacking the appropriate mutational status. Consequently, we developed our own CTC approach, combining a marker dependent and marker independent detection method. We combined positive selection, using CSPG4 and CD146 MACS microbeads and Parsortix TM , in order to prevent the loss of markernegative tumor cells. However, one limitation of our study is the focus on RAS/RAF mutations. Even though RAS/RAF mutated tumors present the majority of mutated melanomas, further research will
Discussion
Our work showed that surface marker expression on melanoma cells is dependent on their mutational status. Furthermore, we demonstrated that a combined analysis of ctDNA and CTCs predicted relapse earlier than imaging, and was more accurate than serum LDH or S100 in a subset of patients. Interestingly, we were able to detect "private" mutations on CTCs and ctDNA, that were not revealed in the random bulk analysis of the primary tumor.
In the present study, we have analyzed melanoma-associated cell surface markers in relation to the mutational status of the melanoma cells. We found a larger proportion of surface marker-positive cells (e.g., CSPG4/MCAM) in the RAS/RAF-mutated cohort compared to the non-RAF/RAS mutated cohort. Thus, we conclude that the commonly employed enrichment of CTCs based on surface marker expression might be biased and could lead to the loss of subsets of tumor cells lacking the appropriate mutational status. Consequently, we developed our own CTC approach, combining a marker dependent and marker independent detection method. We combined positive selection, using CSPG4 and CD146 MACS microbeads and Parsortix™, in order to prevent the loss of marker-negative tumor cells.
However, one limitation of our study is the focus on RAS/RAF mutations. Even though RAS/RAF mutated tumors present the majority of mutated melanomas, further research will have to be conducted to test whether alternative driver mutations might also be represented by specific marker expression.
Overall, 32% of patients were CTC-positive. An increase in CTC-positive patients was detected with increased tumor staging. Enrichment of melanoma CTCs was challenging, due to intra-patient heterogeneity and inter-patient heterogeneity, including different disease stages, subtypes and therapy regimes, as reflected in our patient characteristics ( Figure 3A).
It was previously suggested that ctDNA is more accurate in predicting response to targeted therapy and immunotherapy than serum LDH [24,25]. An increased quantity of ctDNA can be found in the circulation of cancer patients [26]. ctDNA is released from tumor cells via different mechanisms, such as apoptosis, necrosis and secretion [15,26]. The most common mutation in melanoma BRAF can be detected in the ctDNA of melanoma patients and has been shown to be useful in monitoring patients [27]. Sensitive technical strategies for ctDNA detection include ddPCR and BEAMing [6,15,25,28]. Here, we have used a rapid and cost-effective approach for ctDNA analysis, which is based on mass spectrometry, can be used for sensitive multiplex analyses, and requires no bioinformatics. To our knowledge, there is only one previous report using this approach for ctDNA detection in melanoma patients [29]. The panel achieved 92% concordance with ddPCR for the detection of BRAFV600E in ctDNA and was capable of measuring increased levels of mutation in metastatic melanoma patients undergoing therapy prior to radiological progression [29].
Finally, we established a dual approach to detect ctDNA and CTCs and showed proof-of-principle data on two index patients. Patient 2 partially responded to anti-PD1 treatment with Pembrolizumab for 8 weeks and developed pancreatitis; treatment was then discontinued, and the patient relapsed. Intriguingly, CTCs obtained at the time of relapse reveal both a BRAFV600E and EGFRI491M mutation, suggesting a potential benefit from targeted therapy. ctDNA < 150 bp was not detected at this time point, but later, in association with more severe disease progression. Thus, the combined assessment of CTCs and ctDNA can provide complementary information. Patient 3 revealed both CTCs and ctDNA < 150 bp, even during a period of complete clinical remission, in response to BRAF/MEK inhibition. A BRAF and MAPK activating mutation positive CTC was detected during the time of BRAF inhibitor treatment, possibly the first indication of an emerging resistance.
Serum proteins have been frequently used as biomarkers in melanoma in the past. LDH is the only blood-based biomarker implemented in the AJCC melanoma staging system, since the elevated serum LDH is associated with significantly decreased survival in patients with stage IV disease [30]. Nonetheless, LDH is not specific to melanoma or other malignancies; LDH activity can, for example, increase in response to tissue injury of the liver or heart [31]. In addition, levels of serum S100B can indicate clinical response to treatment [32,33]. However, S100 proteins also show an elevated expression in cardiovascular, neurological and inflammatory diseases [33]. Thus, the interpretation of therapy responses using LDH and S100 can be limited, which is reflected in our data. When Patient 2 relapsed at week 18, five CTCs were detected. Neither ctDNA, LDH or S100 levels were elevated. For Patient 3, ctNDA < 150 bp was elevated at week 31, when a DPH3 mutation was detected, possibly a first indication of relapse, which occurred at week 118. At that time point, LDH or S100 levels were within reference values. The relevance of DPH3 mutations in the process of carcinogenesis remains to be determined [34,35]. However, it is noteworthy that DPH3 over-expression was shown to promote cellular invasion and metastasis in murine melanoma cells in vivo, whereas silencing of DPH3 reduced development of metastasis [36].
CTC count and ctDNA < 150 bp appear to be promising tools to predict tumor burden in our cohort. Kaplan-Meier curves demonstrated that patients with detectable CTCs, ctDNA < 150 bp and LDH positive patients show a worse outcome than composite marker-negative patients. This finding is in line with previous studies, where ctDNA levels provide an accurate prediction of tumor response and overall survival in patients treated with PD1 inhibitors [37]. Additionally, baseline ctDNA levels have been found by another group to be significantly associated with progression-free survival in patients treated with BRAF inhibitor therapy [15], and CMCs have shown prognostic value concerning survival in previous studies [38][39][40].
According to the European Society for Medical Oncology (ESMO) guidelines for melanoma, mutation testing on biopsies for treatable mutations is mandatory in patients with advanced disease, to select the appropriate systemic therapy. In cases of inaccessible metastases, liquid biopsy might become a potential approach to guide therapy decisions. The initial mutational status (i.e., mutations in BRAF, NRAS) of the primary tumor was recovered in CTC and ctDNA in 47.6-70.6% of all cases. Importantly, private mutations, not detectable on the primary tumor of the same patient mutations, were detected on CTCs and ctDNA in 29.4-52.4% of all samples, suggesting that liquid biopsy can provide complementary information to analysis of tissue biopsies. Previous studies focusing on BRAF mutations found a concordance between plasma ctDNA and tumor BRAF mutations of 75-76% [41,42]. The detection of mutations, which are not present in the primary tumor, might help to assess tumor heterogeneity and track clonal tumor evolution in individual patients.
Patient Samples
A total of 100 patients were recruited from January 2014 until November 2016 at the Clinic for Dermatology, University Hospital Hamburg-Eppendorf and the Clinic for Dermatology, Elbe-Klinikum-Buxtehude. A total of 84 patients with malignant melanoma fulfilled the inclusion criteria (written informed consent, blood draw, stage I-IV). Patients were staged according to the TNM classification for malignant melanoma (AJCC 2009). Patients of all stages, aged 21-88 years, with cutaneous, uveal, acral and melanoma of unknown primary were included. Blood samples obtained from healthy donors served as a negative control. Blood was drawn into ethylene diamine tetra-acetic acid (EDTA) tubes. The number of CTCs was determined per EDTA-tube (approx. 7.5 mL) of peripheral blood. Written informed consent was obtained from all participants prior to the blood draw, in accordance with the principles and patient rights laid down in the declaration of Helsinki. All laboratory procedures have been approved by the Ethics Committee Hamburg (ethics application PV3779). Our study adheres to the REMARK criteria [43]. Lactate dehydrogenase (LDH) and S100B levels were measured independently by the Department of Pathology, University Hospital Hamburg-Eppendorf.
Tumor Cell Enrichment
To identify the most suitable method to detect different CTC subpopulations, 25 BRAF V600E positive SKMEL28 cells and 25 NRAS/BRAF wildtype MeWo cells (kindly provided by Prof. Dr. med. Udo Schumacher, UKE, Germany) were both spiked into 7.5 mL of blood from healthy donors. Both cell lines were purchased via ATCC, and continuously monitored every three months by STR Profiling and mycoplasma testing (PCR). Subsequently, marker dependent (Cellsearch ® , MACS) and independent (Leucosep TM , Parsortix TM , MACS) [44,45] enrichment methods were tested.
The cells were fixed by adding 700 µL 0.5% paraformaldehyde solution and centrifuged for 10 min at 300× g. After resuspension in 300 µL MACS Buffer (Miltenyi Biotec, Bergisch Gladbach, Germany), the tumor cells were magnetically labelled by adding 20 µL anti-CD146 MicroBead Kit (CD146 MicroBeads and FcR Blocking Reagent, Miltenyi Biotec, Bergisch Gladbach, Germany) and 20 µL Anti-Melanoma (CSPG4) MicroBeads (Miltenyi Biotec, Bergisch Gladbach, Germany) to the cell pellet and incubated at 4 • C for 30 min. After centrifugation (12 min 300× g), 1 mL MACS Buffer was added and the cell suspension was inserted into a MACS separation column (Miltenyi Biotec, Bergisch Gladbach, Germany) that had been equilibrated with MACS buffer. Magnetically labelled cells adhered to the column, while the unlabeled cells passed through. The MACS column was removed from the magnetic field. The labelled tumor cells were flushed from the column with 3 mL MACS Buffer, and additional 3 mL by force. Finally, the cell suspension was centrifuged for 4 min at 1200× g in order to secure cells on a glass slide.
Leucosep ®
Peripheral blood samples were collected in EDTA tubes. After performing plasma separation of the whole blood sample (described below) density gradient centrifugation with Leucosep™ tubes and Ficoll-Paque™ media was used to isolate the peripheral blood mononuclear cells (PBMCs) (800× g, 10 min). The mononuclear cell fraction was transferred to a new 50 mL tube, washed once and centrifuged for 15 min at 300× g, in order to form a cell pellet. After resuspending the cells in PBS, they were transferred to glass slides by cytospin centrifugation.
Negative selection/MACS(−)
Mononuclear cells were prepared as described above (positive selection). CD45 positive cells were depleted from the sample, using anti-CD45 magnetic beads, according to the manufacturer's protocol (Miltenyi, Bergisch Gladbach, Germany).
Marker-independent CTC enrichment (Parsortix TM device)
Parsortix™ is a size and deformability-based method that allows for marker-independent CTC enrichment (Angle Plc, Guilford, UK). Cells were separated according to their size and deformability (final separation gap 8 µm), using a disposable cassette, according to our previous work [45].
Combined Approach (MACS and Partsortix ® )
After adhesion of the magnetically labelled melanoma cells to the column, the column (LS) was washed with 3 mL MACS Buffer (Miltenyi, Bergisch Gladbach, Germany) and the flow-through (4 mL, marker negative melanoma cells and other mononuclear cells) was collected in a Parsortix™ tube and subsequently enriched by the Parsortix™ method. For the isolation of marker positive cells, the MACS column (Miltenyi, Bergisch Gladbach, Germany) was removed from the magnetic field. The labelled tumor cells were flushed from the column with a 6 mL MACS Buffer.
Immunofluorescence Staining
After enrichment, cells were transferred to the cytospin (max. 1 million per slide, 3 min 1200 r.p.m.) and dried overnight. After fixation with 0.5% paraformaldehyde solution for 10 min, cells were stained for surface markers CSPG4 (anti hNG2/MCSP, R&D Systems, Minneapolis, MN, USA) and MCAM (anti-CD146 monoclonal antibody, Merck KGaA Darmstadt, Germany) (positive markers) or cytoplasmic S100 (Anti-S100, Dako Denmark A/S, Glostrup, Denmark) (to enable detection of surface marker negative cells), the common leukocyte antigen CD45 (FITC anti-human CD45, BioLegend R , San Diego, CA, USA) (negative marker), and the nuclear dye DAPI. The slide was incubated with the respective antibodies for 1 h at room temperature, or overnight at 4 • C. Enriched cells were quantified by fluorescence microscopy. Morphologically intact NG2+/MCAM+/CD45−/DAPI+ cells were defined as CTCs and picked with a micromanipulator. Single cells were stored at −80 • C for future amplification and mutational analysis.
Whole Genome Amplification
Whole genome amplification (WGA) on isolated CTCs was performed using the Ampli1 Kit (Silicon Biosystems, Castel Maggiore, Italy), according to the manufacturer's instructions. Quality of the WGA product was analyzed using the Ampli1 QC Kit (Silicon Biosystems, Castel Maggiore, Italy).
cfDNA Extraction
Blood samples were collected in EDTA tubes, stored at RT and processed within 6 h. Shipped blood samples were stored in StreckTubes and processed within 36 h. In order to isolate the plasma from the whole blood, the samples were double centrifuged for 10 min at 300× g. Plasma was transferred to a new tube and centrifuged at 2000× g for 15 min to remove cellular debris. Plasma aliquots were stored at −20 • C/−80 • C. Cell-free DNA (cfDNA) was isolated from 1-5 mL plasma with the QIAamp Circulating Nucleic Acid Kit (Qiagen, Hilden, Germany), according to manufacturer's instructions, with a final elution volume of 40 µL.
Quantification and Size Fragment Distribution of cfDNA
The concentration of cfDNA was determined using a NanoDrop Spectrometer ND-1000 (Thermo Fisher Scientific, Waltham, MA, USA) with a sample volume of 1 µL. Fragment distribution was assessed using the 4200 TapeStation device, using the High Sensitivity D5000 ScreenTape Assay with 1 µL sample, and 5 µL High Sensitivity D5000 Sample Buffer (Agilent, Santa Clara, CA, USA).
Mutational Analysis
Mutational analysis was performed using the UltraSEEK™ Melanoma Panel v1.0 (Agena Bioscience, Hamburg, Germany), interrogating 61 clinically relevant variants across 13 genes, including BRAF, NRAS, KIT and MAP2K1, detected at as low as 0.1% minor allele frequency. Reactions were performed as described before [46]. In brief, PCR (45 cycles) was followed by shrimp alkaline phosphatase treatment and single base primer extension, using biotinylated ddNTPs specific for the mutant alleles. After capture of the extended primers using streptavidin-coated magnetic beads, a cation-exchange resin was added for cleaning, and 10-15 nL of the reaction was transferred to a SpectroCHIP ® Array (a silicon chip with pre-spotted matrix crystals) using an RS1000 Nanodispenser (Agena Bioscience). Data were acquired via matrix-assisted laser desorption/ionization time-of-flight mass spectrometry, using a MassARRAY Analyzer 4 (Agena Bioscience, Hamburg, Germany). After data processing, a spectrum was produced with relative intensity on the y-axis and mass/charge on the x-axis. Typer Analyzer software was used for data analysis and automated report generation. Sanger sequencing was performed to verify mutations detected by the UltraSEEK™ Melanoma Panel, and only mutations which were detected in both assays (98%) were used for further analysis.
RNA Analysis
CTCs were isolated within 4 h of blood withdrawal. cDNA synthesis and amplification were performed using the SuperScript II Kit (Thermo Fisher Scientific, Waltham, MA, USA), according to the manufacturer's recommendations.
Bioinformatical and Statistical Analysis
TCGA data was last accessed via the following webpage in November 2017-http://firebrowse.org/. Differential analysis was performed using the R packages edgR and Limma. Go-term analysis of genes upregulated in the RAS/RAF group FC > 2 adj. p-value < 0.05 was performed using the ClueGo app in Cytoscape. Single cell analysis of the Tirosh et al. dataset (GSE72056) was performed using the R packages SingleCellExperiment. Data were plotted using ggplot2.
The distribution of disease specific survival was estimated using the method of Kaplan-Meier. Median values for the cox regression model for the distributions of the HRs and P values are reported with 95% empirical. Analysis was performed using the R-packages survminer.
Statistical analysis was performed using GraphPad Prism Software (GraphPad Software Inc., La Jolla, CA, USA). All datasets are represented as mean ± SEM and were analysed either by ANOVA and Tukey's or Holm-Sidak's multiple comparison correction. Statistical significance was considered at p-values of p < 0.05.
Conclusions
In summary, analysis of CTCs in combination with ctDNA provides complementary information, beyond the current serum biomarkers LDH and S100, which might help to personalize targeted and immunotherapies for melanoma patients in the future. However, the present findings need to be validated in larger future studies before implementation into clinical practice. | 9,413.6 | 2019-10-29T00:00:00.000 | [
"Medicine",
"Biology"
] |
Microwave spectroscopic study of the hyperfine structure of antiprotonic helium-3
In this work we describe the latest results for the measurements of the hyperfine structure of antiprotonic helium-3. Two out of four measurable super-super-hyperfine SSHF transition lines of the (n,L)=(36,34) state of antiprotonic helium-3 were observed. The measured frequencies of the individual transitions are 11.12548(08) GHz and 11.15793(13) GHz, with an increased precision of about 43% and 25% respectively compared to our first measurements with antiprotonic helium-3 [S. Friedreich et al., Phys. Lett. B 700 (2011) 1--6]. They are less than 0.5 MHz higher with respect to the most recent theoretical values, still within their estimated errors. Although the experimental uncertainty for the difference of 0.03245(15) GHz between these frequencies is large as compared to that of theory, its measured value also agrees with theoretical calculations. The rates for collisions between antiprotonic helium and helium atoms have been assessed through comparison with simulations, resulting in an elastic collision rate of gamma_e = 3.41 +- 0.62 MHz and an inelastic collision rate of gamma_i = 0.51 +- 0.07 MHz.
super-hyperfine (SHF) splitting, which can be characterized by the angular momentum Sp. Even though the magnetic moment of the antiproton is larger than that of the 3 He 30 nucleus, the former has a smaller overlap with the electron cloud. Therefore it creates 31 a smaller splitting. The complete hyperfine structure for p 3 He + is illustrated in Fig. 1.
3. Laser-microwave-laser spectroscopy 50 The first observation of a hyperfine structure in antiprotonic helium was achieved in subsequently Auger decay of the transferred atoms and annihilation of the antipro-71 tons in the nucleus will occur. The number of annihilations after the second laser 72 pulse will be the larger the more antiprotons were transferred by the microwave pulse.
74
When the antiprotons first enter the helium gas, a large annihilation peak ("prompt 75 peak") is caused by the majority of formed pHe + atoms which find themselves in Auger 76 decay-dominated states and annihilate within picoseconds after formation. At later 77 times, this peak exhibits an exponential tail due to pHe + atoms in the metastable Since the intensity of the antiproton pulse fluctuates from shot to shot, the peaks 87 must be normalised by the total intensity of the pulse (total). This ratio is referred to 88 as peak-to-total. The peak-to-total (ptt) corresponds to the ratio of the peak area (I(t 1 ) 89 or I(t 2 )) to the total area under the full spectrum. If the second laser annihilation between antiprotonic helium atoms and regular helium atoms. Refilling from higher-98 lying states also contributes to the equalization of the hyperfine substate populations.
99
In general, a short delay T is preferable because the signal height will decrease for 100 longer laser delay times as a result of the exponential decay of the metastable state 101 populations. However, the linewidth of the RF transition will increase if the delay is too 102 short. Further, far higher RF power will be required to complete one spin-flip. If the 103 delay is too long, the collisional relaxation of the system would already have eliminated 104 any asymmetry between the two states caused by the first laser pulse. The signal would 105 be too low to be observed.
106
The two pulsed lasers were fixed to a wavelength of 723.877 nm, with a pulse titanium window for the antiproton beam and a 4 mm thick fused silica window for the 135 laser beam to enter [23], and are equipped with meshes to contain the microwaves.
136
In order to measure the annihilation decay products two Cherenkov counters are 137 mounted around the target volume, connected to photomultipliers (cf. Fig. 3) . They
157
In preparation for the actual investigation of the hyperfine substructure, via microwave 158 resonance, several studies are required to optimize the parameters such as laser power, 159 laser resonance frequency, laser delay time and microwave power. To fit the two transitions, a function of the natural line shape for a two-level system 210 which is affected by an oscillating magnetic field for a time T was used. It is given by [26] 211 Here X(ω) is the probability that an atom is transferred from one HF state to the other, HF for the fitting of the raw data together with the reduced χ 2 /ndf and ν HF after inflating the errors of the individual data points by χ 2 /ndf . The fit transition frequencies are displayed for the two different fitting methods, ASF and ISF. At the higher resonance the frequency points differed slightly between 2010 and 2011. These data can only be combined in the averaging over all single scans. The microwave power for the 11.157 GHz resonance was further lower by about 2.5 W compared to 2011. Therefore, the values obtained by the ISF method were used as final results. HF and ν −+ HF in comparison with three-body QED calculations, where ν HF denote the SSHF transition frequencies, δ exp is the relative error of the measured frequencies and Γ the resonance line width. The relative deviation of experiment and theory is defined as δ th−exp = (ν exp − ν th )/ν exp . The quoted theoretical precision is ∼ 5 × 10 −5 from the limitation of the Breit-Pauli approximation that neglects terms of relative order α 2 . This does not include numerical errors from the different variational methods used. For ref. [11] such as microwave power, Q value and laser delay, the measured values were taken.
257
To assess the rates of collisional effects which induce relaxations between the SSHF Table 2), given as linear frequencies, confirmed that the density dependence is very small. Also for p 3 He + theory predicts a 314 collisional shift at the kHz level, much smaller than the experimental error bars [17].
315
For the frequency difference ∆ν ± HF = ν −+ HF −ν −− HF between the two SSHF lines around 316 11 GHz there is an agreement between both theoretical results and experiment within
324
The two transitions at 16 GHz could not be measured anymore due to lack 325 of beamtime -even though the microwave target was readily tested and calibrated.
326
However, we came to the conclusion that the observation of these two resonance lines 327 would deliver no additional information on the investigated three-body system and 328 primarily serve to accomplish a complete measurement of the p 3 He + hyperfine structure.
329
This study with p 3 He + was considered a test of QED calculations using a more | 1,542.8 | 2013-03-12T00:00:00.000 | [
"Physics"
] |
A curve-based material recognition method in MeV dual-energy X-ray imaging system
High-energy dual-energy X-ray digital radiography imaging is mainly used in the material recognition of cargo inspection. We introduce the development history and principle of the technology and describe the data process flow of our system. The system corrects original data to get a dual-energy transparence image. Material categories of all points in the image are identified by the classification curve, which is related to the X-ray energy spectrum. For the calibration of classification curve, our strategy involves a basic curve calibration and a real-time correction devoted to enhancing the classification accuracy. Image segmentation and denoising methods are applied to smooth the image. The image contains more information after colorization. Some results show that our methods achieve the desired effect.
Introduction
The X-ray imaging technique has become one of the most important tools in customs inspection. Presently, there are mainly two X-ray imaging modalities: radiography and computed tomography (CT). Although CT can provide 3-D structures and an accurate attenuation map of the cargo, its complexity and high price limit its application [1][2][3]. X-ray radiography, including single energy and dual energy, is still the mainstream technology. The development of X-ray radiography undergoes three stages: X-ray film photography, computed radiography (CR) and digital radiography (DR). The single-energy X-ray DR image merely gives the cumulative density information of the irradiated objects in one direction. It is used in preliminary medical diagnosis and simple security inspection. Since single-energy X-ray DR provides limited information, the dual-energy method was developed. Low-energy dual-energy X-ray DR imaging has been widely used in current security inspection equipment, which can detect and distinguish contraband by determining material atomic number Z. The X-ray's energy here is usually lower than 1 MeV. This technology is inapplicable in high Z material recognition or cargo inspection, as the energy of the X-ray which can penetrate the object in these situations needs to be a few MeV.
The British company Cambridge Imaging first proposed the idea of high-energy dual-energy X-ray imaging. There were some disputes about the validity of the high-energy dual-energy method in material recognition. The Russian Efremov Research Institute proved the feasibility of this method with their experimental prototype [4]. The German company Heimann and the American company EG&G applied X-ray hardening technology to this field and proposed the filter method. The Department of Engineering Physics at Tsinghua University and its cooperative enterprise, Nuctech, established a platform and made some achievements on material recognition and related studies. The theory of high-energy dual-energy X-ray DR imaging and material recognition has been deeply studied, and the corresponding experiment results further validated the feasibility of dual-energy imaging material recognition [5].
In this paper, we construct an imaging system model and a whole data processing flow. For the best visual effects of the final results, we used the image smoothing strategy and image colorization processing. Some realization details are also given. The R-curve material recognition method is a typical high-energy dual-energy X-ray DR material recognition method [6]. We developed a real-time R-curve calibration method. It deals with the differences of the Rcurves of different energy spectra caused by the system status fluctuation and inconsistency. In Sect. 2, we introduced the technology principle and elaborated on the methods of a MeV dual-energy imaging model. We focused on the calibration strategy we designed. In Sect. 3, we gave and discussed some experimental results. We concluded and envisioned future work in Sect. 4.
2 Theory and method 2.1 The principle of MeV dual-energy X-ray imaging in material recognition The three main interactions between a photon and matter are the photoelectric effect, Compton scatter effect and the electron pair effect. They, respectively, dominate the low (\1 MeV)-, middle (1-3 MeV)-and high ([3 MeV)-energy range [7]. The corresponding attenuation coefficients, l, have different dependences with a material atomic number, Z. We can give where P, CS and EP are the abbreviations of the three effects. Consider an X-ray source whose energy spectrum is N(E) and the highest energy is E m , a single substance with an atomic number of Z, an attenuation coefficient function of lðE; ZÞ and a thickness t, the transparence, T, is In a dual-energy situation, the boundary energy of the two X-ray sources is E 1 and E 2 . We defined the logarithmic ratio of T as where R is the ratio of the equivalent attenuation coefficient, l. When the X-ray source is monochromatic, which means the energy spectrum N(E) is a single line, Eqs. (2) and (3) can be simplified. Suppose E 1 is in the low-energy range and E 2 is in the middle energy range, R can be written as From Eq. (3), R is easily computable. From Eq. (4), R is a clear indication of Z. Besides, when the X-ray is polychromatic, a great dependence between R and Z still exists. Based on these facts, low-energy dual-energy X-ray DR imaging technology has been widely used in small security inspection devices for material discrimination. Low-energy dual-energy means that E 1 and E 2 are usually lower than 1 MeV. When Z is high and the irradiated object is thick, lowenergy X-ray imaging becomes useless. If we change E 1 to the high-energy range and keep E 2 in the middle energy range, then Eq. (4) is and R is dependent on Z to a certain extent, so it can be still used to classify material. This ideal conclusion without consideration of the subordinated interaction of photons with matter is based on the assumption of the single-line Xray energy spectrum. In fact, a MeV X-ray DR system uses the linear accelerator as the X-ray source, which generates X-rays with a broad energy spectrum [8]. Most of the photons were distributed in the middle energy range, where l and Z have no correlations. The effectiveness of R in material recognition is not obvious. It was found that, although R changes when thickness t changes, R is still dependent on Z [9]. Here E 1 and E 2 are usually higher than 3 MeV.
A MeV dual-energy system
We use a schematic model (see Fig. 1), including an accelerator, which emits a vertical fan-shaped X-ray beam, a scanning track, which is perpendicular to the X-ray's main beam direction, an L-shape detector and a data processing unit. The cargo moves along the scanning track, while the L-shape detector receives the photons passing through the cargo and form the dual-energy X-ray images.
The data processing unit consists of four steps. First, correct the acquired original dual-energy X-ray images and calculate the dual-energy transparence images. Second, use the classification curve on the dual-energy transparence images to form the material information image. Third, to improve the image quality, a smoothing process is implemented. The final step is the colorization of the gray image. In the next four sections, we will introduce these four steps of the data processing unit and mainly concentrate on the calibration of the classification curve. We propose a real-time R-curve calibration method. In Sect. 3, we can see that our method enhances the classification accuracy and gives a satisfactory visual result.
Data acquisition and preprocessing
We assume that the accelerator produces the dual-energy X-ray simultaneously. Accordingly, the detector is able to distinguish high and low-energy X-rays and form the dual-energy X-ray images separately. Other aspects of our model are basically the same as reality. The fan-shaped X-ray has an angular distribution. Its main beam is located near the middle of the fan and has a maximum intensity decreasing toward both sides. It causes different vertical positions of the X-ray image which have a different X-ray intensity and energy spectrum. The accelerator state fluctuation in the scanning process causes different lateral positions of the X-ray image which have a different X-ray intensity and energy spectrum. The detector background and response inconsistencies also exist. In data processing unit, the first step is to correct the original data and get the dual-energy transparence image. We give Eq. (6) based on Eq. (2).
The coordinate x, y represents the lateral and vertical position. I is the original dual-energy image. I 0 is the air image. I BK is the detector background image. The correction factor, LD(x), obtained by monitoring the accelerator state fluctuation in the scanning process is a function of the lateral position, x. In Eq. (6), the division by I 0 corrects the intensity angular distribution and detector response inconsistencies. Rest corrections are also done by Eq. (6). After the point-by-point correction and simple denoising, we can get the dual-energy transparence image, T.
To take advantages of two transparence images, we use a pointwise weighted fusion of them. In the thick places of the irradiated object, the high-energy transparence image gives more information than the low-energy transparence image. The conclusion is opposite in the thin places. The gray value, representing one point's thickness, determines the weight. Using the dual-energy transparence image and the classification curve, which we will elaborate on next, we can get the material information image. The colorization assembling the fused transparence image and the material information image gives a final result.
Curve-based material recognition method
Four kinds of classification curves, R-curve (Fig. 2a), T 2 À T 1 curve (Fig. 2b), a-curve (Fig. 2c) [10,11], and a separability curve which denotes the transparence difference of the two materials (Fig. 2d), are shown in Fig. 2. Among them, the R-curve has the best visual separability, so we use the R-curve to show the calibration output and result comparison. The data in Fig. 2 are ideal. The longitudinal coordinate of the R-curve is R, defined in Eq. (3), and the horizontal coordinate is the inverse of T 1 or T 2 . Different materials have different R-curves. They are arranged in Z's increasing order when they are in the same image [6]. The classification curve here has several single R-curves related to the same number of typical objects. We assume that there are four typical objects, including Pb, Fe, Al and CH 2 (or C), as we can see in Fig. 2.
(x, y) is one point of the dual-energy transparence image. There are two transparence values, T 1 ðx; yÞ and T 2 ðx; yÞ. They give a point, ðR; 1=T 2 Þ, on the classification curve, C. Because R-curves of different materials are arranged in Z's increasing order, the point (x, y) is obtained by interpolating between two adjacent R-curves in C. By repeating this procedure on each point of the dual-energy transparence image, the material information image Z(x, y) is formed.
Equations (2) and (3) show that the R-curve correlates with the energy spectrum. We already know that each point of the dual-energy transparence image has different X-ray intensity and energy spectrum because of the angular distribution and the accelerator state fluctuation. This fact causes a problem that different points of the dual-energy transparence image need different classification curves. However, it is impossible to calibrate all of the classification curves. We designed a new calibration method to obtain all of the points' approximate Rcurves.
The calibration strategy we employ takes two steps. First, we get the basic classification curve before the system scans cargo or after the system state significantly changes. This step requires scanning several typical materials. In one scanning process, the system only scans one typical material with a different mass thickness. Like the common scan, we get two transparence images T 1 , T 2 , and each point, (x, y), belongs to an R-curve, R xy . Assuming a steady system state in the process, we ignore the differences of energy spectrum in different lateral positions and just use several vertical points to represent the whole angular distribution. So we have The data are data y ¼ fR y ½T 1 ðx i ; yÞ; T 2 ðx i ; yÞ; i ¼ 1; 2; . . .; mg; where m is the number of different mass thicknesses and n is the number of vertical points. The R-curve is fitted with these data, and the classification curve is formed with these fitted curves of several typical objects. Here, we use Pb, Fe, Al and C.
fR Z y jR Z y ¼ fitðdata Z y Þ; y ¼ y 1 ; y 2 ; . . .; y n g; C y i ¼ fR y i ;Z jZ ¼ Pb; Fe; Al; Cg; i ¼ 1; 2; . . .; n: In the cargo scan, the system monitors the state variation to complete the real-time calibration of the classification curve. A small device consisted of the typical materials with the single thickness set in a certain vertical position, Y, and scanned synchronously with the cargo. The data are where nx is the number of horizontal pixels. In the first step, these data are also saved as T Z 1b ðx i ; YÞ; T Z 2b ðx i ; YÞ. They are the average of T on the horizontal direction, as we ignore the lateral difference. Then the revised classification curve will be C x j y i ¼ fR Z x j y j jR Z x j y j ¼ R Z y i à FðT Z 1r ; T Z 2r ; T Z 1b ; T Z 2b ; x j ; YÞ; Z ¼ Pb; Fe; Al; Cg; i ¼ 1; 2; . . .; n; j ¼ 1; 2; . . .; nx: The function F uses real-time calibrating data to estimate the difference between R Z xv and R Z v . The correction method in Eq. (11) relies on how function F is calculated. Considering the statistical straggling of Tðx j ; YÞ, we choose a segment of x j and use the average of the monitoring data to correct the classification curve.
Smooth of material information image
The detector signal noise inevitably exists, and it is even amplified in the data processing. It has an impact on material recognition accuracy and makes the material information image rough. The visual effect of final result after colorization is not good enough. Material information image quality promotion is necessary. Some literature proposes image segmentation smooth strategy [9,12]. We apply the same idea.
The first step of our smooth process is the segmentation of the fused transparence image, which is obtained previously. The image is segmented into regions, which keep the continuity of the irradiated objects interior as much as possible and discriminate different irradiated objects as clearly as possible. The average of all the Z values in a corresponding region in the material information image is assigned to all points in this region.
The general image segmentation algorithm, like the single-pass split-merge algorithm, or data clustering algorithm, like the Leader algorithm, can be used here with some adjustment [13]. The irradiated objects may be mixed and disordered. So the segmentation or cluster result may have too many small areas with only several pixels. This over-segmentation can be solved by merging the small area into the nearest large area. The 'nearest' means not only the distance but also the similarity between them.
Using the average of the Z values in one segmentation region to replace all points in this region brings some loss of the original information. The denoising approaches can give a better material information image, and the majority of the image remains the same. It is a challenge to find a better method, which can smooth the image while maximizing the retention of the original information.
Colorization
The idea of colorization and the IHS color space was proposed in Ref. [9,12]. We use a similar approach applying HLS color space [14]. In the colorization, different colors represent different materials. We use three color spaces, including RGB, HLS and YUV. If all three values in a color space are known, a color is determined. We divide the range of Z into several parts as follows: There are p þ 1 hues H 1 , H 2 , . . ., H pþ1 . When the Z value of one point in the material information image falls into the jth part, the H value in the HLS color space equals H j . The sensitivity of the human eye to different color is different. There are red, green and blue with the same L value in the HLS color space. The brightness felt by eyes is different. Green is the brightest and red is brighter than blue. If the L value equals the gray value, the points with same gray value will give a different brightness when we look at the final result. It is inappropriate.
In the YUV color space, colors with same Y value give the closest brightness feeling. Let the Y value equal the gray value of the fused transparence image. The YUV color space and HLS color space can be converted to each other. So we have As Y, H, f is known and S is given, L is the solution of Eq. (13). Then all three values in the HLS color space are already set, and a color is determined. Repeat this procedure on every point of the material information image to get the final result. The table of hue and the saturation value, S, are changeable and directly influence the image's visual effect. The mapping relationship between Y and the gray value can be optimized by some transformation, like logarithm stretching.
Experimental result
The data are provided by a 6/3 MeV X-ray DR imaging cargo inspection system, which applies our basic design of a data processing unit. The accelerator alternatively emits the high and low-energy X-rays. The emission frequencies are both 40 Hz. We use a single column of the CWO detector and a scanning speed is 0.2 m/s. The basic classification curve is formed by scanning three single materials, C, Al and Fe. A stair-step object with a single material component is scanned to form one R-curve [15]. All R-curves together form a classification curve. For clarity of the results, we use one R-curve representing a classification curve. There are two different system states. In Fig. 3, the dotted curve is the classification curve in state 1 and the dashed curve represents the classification curve in state 2. The solid curve is the revised classification curve based on the classification curve in state 1. The revision uses the difference between the real-time calibrating data of the two states. We think that the dashed curve is the 'true' classification curve in state 2 and the solid curve is an estimation of the 'true' one. Their closeness shows the effectiveness and rationality of our calibration strategy.
We arrange eight objects in the order of Pb, Fe, Al, C, Al, C, Pb and Fe and divide them into two groups according to size. Their mass thicknesses are all 40 g=cm 2 . The scanning of these eight objects is in state 2. Suppose we do not know the classification curve in state 2 (dashed curve in Fig. 3). Our calibration strategy means that we can get the classification curve in state 2 if we know the classification curve in state 1 and the real-time calibration data of the two states. In Fig. 4, the image a) on the left is the final result without the use of the real-time calibrating data. For the four larger objects Al, C, Pb and Fe, the specific Z values obtained using the classification curve in state 1 (dotted curve in Fig. 3) are 46, 27, 62 and 54. The image b) in the middle is the final result with the use of the real-time calibrating data. For the four larger objects Al, C, Pb and Fe, the specific Z values obtained by using the calibrated curve (solid curve in Fig. 3) are 18, 9, 54 and 45. The image c) on the right is the color table. The system colorization settings give that the hues of four typical objects, C, Al, Fe and Pb, are orange, green, blue and purple. We also use the classification curve in state 2 (dashed curve in Fig. 3) and get the Z values 16, 5, 54 and 45. Note that there is no R-curve of Pb, and the Z values of the Pb object from three sets of results are 62, 54 and 54 and far from 82. The consequent color of the Pb object is deviated from the righteous color. Although the Z value of the Fe object has an obvious deviation due to the lack of the R-curve of Pb, it is in the righteous region, and thus, the color is also righteous. The data and figure comparison clearly show the effectiveness of the real-time calibration and also match the comparison of the curves in Fig. 3. Besides Pb, the other materials' results show that our calibration strategy enhances classification accuracy.
Notice that all objects in Fig. 4 are made up of the uniform single material, and the two classification curves under two states have distinct differences. The conditions are ideal, and accordingly, the results are good. In the real scan, the scanned objects are always complex. We may be unable to figure out true Z values of all the points, and thus, we cannot verify the accuracy of the recognition results. What we can do is to make the classification curve, whether directly calibrated or real-time revised, as accurate as possible. There are two points to note. First, the calibration of the basic classification curve is used as the basis for a steady system state. Second, the real-time monitoring data have statistical fluctuation even exceeding the state fluctuation of the accelerator or angular distribution. Using the average of a piece of data to reflect the variation is better than using every single point.
In Fig. 5, the comparison is the final color image with and without the smooth process. The larger orange object in image a) looks not uniform, although it should be the same color. We can see that the smoothing improves the image quality. However, the effect is obsolete because of the monotony and uniformity of the irradiated objects. When the objects are complex, the smoothing process will influence the classification accuracy because the rearrangement of the Fig. 3 (Color online) The 'true' curves in two states are C 1 and C 2 . The revised R-curve using the real-time calibration data to estimate the 'true' curve in state 2 is C m2 In Fig. 6, we give a cargo inspection result. The emission frequencies of the dual-energy X-rays are both 33 Hz. The scanning speed is 0.2 m/s. The irradiated objects from left to right are a cigarette, salt, sugar, coffee, buckets of water and concrete. According to the continuity of the irradiated objects, we can assure that the spots and stripes in the object regions of the image on the top are noise and need to be removed. The smoothing process eliminates the noise spots and nonuniformity and significantly promotes the image quality. In the red circled region of the image on the top, the bottom margin of the bucket is overwhelmed by the noise and almost disappears. After the smoothing process, the margin is recovered. Our smoothing method may strengthen the details of simple object regions.
In the cigarette region, the thickness of the cigarette is small. When the irradiated object is thin, the calibration curve separability is worse. With the data fluctuation, the final color image will be full of stripes and spots, as we can see in the amplified region of the image on the top. The cigarette belongs to the orange category, so the blue pixels are noise. The smoothing process takes the average in the cigarette region of the material information image. So the final color will be the middle color between orange and blue, according to the color table. The color changes to
Discussion and conclusion
We described a simple MeV dual-energy X-ray DR imaging cargo inspection system with a detailed description of the data processing unit. The preliminary treatment converts the original data into images for subsequent processing needs. The calibration strategy of the classification curve enhances the classification performance. The smoothing of the material information image is to enhance the image quality. Segmentation is the key part of our smoothing process. Better segmentation methods lead to better image quality. Color imaging can carry more information and give a better visual effect. Colorization can be regulated in different application environments. This system design has a certain guiding significance to the engineering.
Our calibration method is devoted to give the correct classification curve and enhance the classification accuracy. There are several different directions to achieve this goal or push forward the further development of dual-energy X-ray imaging technology. Add a low-energy detector to promote recognition ability of thin objects [16]. Use the Cerenkov detector, which has a threshold and a good response to high-energy X-rays [17]. Add the obstacle's classification curve on the blocked material's classification to enhance the blocked material's recognition accuracy [18]. Use small angle scatter to realize the material recognition [19]. These methods can be applied to our system or a new imaging system based on our model and data processing flow. | 5,880 | 2016-02-01T00:00:00.000 | [
"Engineering",
"Materials Science",
"Physics"
] |
An Integrated Framework to Achieve Interoperability in Person-Centric Health Management
The need for high-quality out-of-hospital healthcare is a known socioeconomic problem. Exploiting ICT's evolution, ad-hoc telemedicine solutions have been proposed in the past. Integrating such ad-hoc solutions in order to cost-effectively support the entire healthcare cycle is still a research challenge. In order to handle the heterogeneity of relevant information and to overcome the fragmentation of out-of-hospital instrumentation in person-centric healthcare systems, a shared and open source interoperability component can be adopted, which is ontology driven and based on the semantic web data model. The feasibility and the advantages of the proposed approach are demonstrated by presenting the use case of real-time monitoring of patients' health and their environmental context.
Introduction
Worldwide healthcare standards are considered important indicators of human progress and civilization as they strongly affect both the economy of countries and the quality of life of citizens [1]. While citizens' expectations from healthcare systems increase, healthcare is facing a high risk of severe decay in the near future, as demographic changes are causing the rise of radical costs and staff shortages in the healthcare system of many countries. The population is ageing in many geographic areas and is growing in others; the percentage of 65+ population in the EU15 will grow from the current 17.1% to the 25% expected in 2030 and the life expectancy will go up to 81.4 years in 2030 from the 75.9 years of 2008. Concurrently, the number of people with chronic diseases, with comorbidities, or with some kind of impairment is increasing. Chronic diseases represent the greatest cause of early death and disability. Cardiovascular disease (CVD) is the leading cause of death in Europe and in the industrialized countries. Accordingly, healthcare and social costs are exploding from a worldwide average of 9% of the gross domestic product (GDP) in 2005 to 11% expected in 2015. Most EU member states spend around 30%-40% of the total health expenditure on the elderly population and long-term care (as an example, Europe spends around 3% of the GDP for treating CVD) [1].
The Institute for Healthcare Improvement reports that many healthcare systems around the world will become unsustainable by 2015 and the only way to avoid this critical situation is to implement radical changes. As a consequence, the healthcare system is subject to reform in many countries. The challenge is to improve care efficiency and effectiveness and to support sustainability of healthcare evolution. The goal is to manage increasing costs, and reduce unnecessary tests by having access to all relevant data and promoting integration of diagnosis and treatment. The common approach to improve the quality of the care process is to enable service access "at any time and in any place" and to move from "how to treat patients" to "how to keep people healthy and prevent illness." Emphasis is put on prevention instead of treatment, supporting early discharge, less costly recovery, and rehabilitation at home [2].
Information and communications technology (ICT) plays an important role in this change. Particularly, ICT contributes to reduce the distance between the user and the clinician, by enabling telemonitoring, which is based on the integration of healthcare, telecommunications, and information technologies. Telemonitoring enables continuous observation of patients' health status and offers the opportunity to collect and analyse large amounts of data regarding the patients and their clinical history. Different solutions exist and they meet specific requirements concerning security, privacy, usability, robustness, safety, and data traceability.
Already in the year 2000, state-of-the-art analysis of remote care solutions demonstrated the benefits of telemedicine to healthcare costs [3]. Still, there is a recognized barrier to healthcare radical innovation and associated cost-effectiveness improvement: this barrier is the fragmentation of healthcare solutions and the lack of interoperability at many levels [4]. Bringing interoperability into the healthcare system is a great challenge as it requires innovation not only in healthcare technology in general but also in management and working style [5].
This paper focuses on one technological aspect, which is information level interoperability, and claims that the wide adoption of interoperability platforms is the way to open innovation in healthcare. This will be demonstrated by showing how an interoperability platform developed within the framework of a European project [6] can enforce information interoperability between different healthcare "legacy" solutions.
The paper is organized as follows. Section 2 reviews the history and current research in out-of-hospital care and shows the impact of lack of interoperability in health care systems. Section 3 describes the proposed approach and summarizes CHIRON's vision [7]. The proposed interoperability platform, developed within SOFIA project is then described in Section 4. In Section 5, we present a general scenario based on this platform to provide innovative services obtained by the concurrent monitoring of people's physiological parameters together with their surrounding environmental conditions. Section 6 describes our interoperable and cross-domain application, together with its architecture and its design process. In Section 7 the conclusions are drawn.
Motivations and Related Works
The need for telemedicine was already felt back in the seventies [8], but telemedicine started to develop fast only in the nineties, when data communication services over the telephone lines became widely available (e.g. [9]).
By that time the social role of telemedicine in healthcare started to be investigated by public institutions like EU (1990) and World Health Organization (WHO/OMS); (1997) and, consequently, a wide number of telemedicine systems were then realized, the main purpose being to reduce the amount of time a primary physician must spend with the patient and to allow at the same time a high level of care.
At the beginning, these telemedicine systems were dedicated to a specific application (e.g. telecardiology, teleradiology, telespirometry, and teledialisis). Then, with the advent of internet-enabled platforms, including mobile platforms, and with the development of location-based services, remote collection of contextualized data from sensors (on-body, stationary, located at home, or in other environments) became possible. Sensors and also innovative actuators opened up new healthcare scenarios and new opportunities in e-health services, many of which are investigated by publicly cofunded research projects.
Classical telemedicine projects are focused on dedicated solutions. Their target is mostly limited to sensor data collection and delivery to medical personnel. Data mining and decision making are marginally addressed. These kinds of systems are usually based on specific sensors. It is difficult to add devices from different manufacturers, with different protocols or based on different technologies. For example, OLDES (Older people's e-services at home) [10] is an EU cofunded project under the IST Programme aiming to plan and develop a technological solution for teleassistance and teleaccompany; it supports a predefined set of specific instruments and VoIP technology. A similar project is AILISA [11], a French initiative promoting an experimental platform to evaluate monitoring technologies for elderly people in their homes. All of these projects benefit from the progress of ICT technologies; they are addressed to the elderly people and implement the classical telemedicine approach, and they are not able to exploit all the potential offered by such technologies because they do not inherently provide an automatic decision flow nor do they offer interoperable solutions.
These features can be found in "health smart homes" targeting elderly and impaired people at home. Smart homes are presented as a "dwelling incorporating a communications network that connects the key electrical appliances and services and allows them to be remotely controlled, monitored, or accessed" [12]. In particular a health smart home should ensure "an autonomous life at home to people who would normally be placed in institutions" [13]. These solutions are based on ambient intelligence research results and try to adapt the technology to people's needs by building on three basic concepts: ubiquitous computing, ubiquitous communication, and intelligent user interfaces. A proof of the importance of the integration of context data into remote health care applications is the growing number of projects and systems with this purpose. In this respect, we can cite MonAMI [14], Smart Medical Home [15], Aware Home [16], TAFETA [17], Gator Tech House [18], Duke Smart Home [19] and B-Live system [20]. Other solutions are discussed in [21]; this survey presents an international selection of advanced smart home solutions.
The common denominator within the above-mentioned solutions is that every system maintains its own proprietary protocols, platforms, information store, and data representation models. Therefore the information produced remains inside the original system implementation leading to an inherent fragmentation along the health care cycle.
This fragmentation, due to the lack of an ontology or a common standard, is the source of unnecessary overhead and it is also a barrier to healthcare service innovation.
If health and context data are shared and made interoperable, then they might be combined and the new knowledge thus generated could support innovative applications for the benefit of multiple institutions and users. Interoperability is International Journal of Telemedicine and Applications 3 becoming a more and more relevant requirement for which solutions have been proposed, for example, by the database community [22] even before the start of the smart houses research era. Interoperability between heterogeneous and distributed systems handling information originated by the physical space has been increasingly investigated since the emergence of Weiser vision on pervasive computing [23]. The IEEE defines interoperability as "the ability of two or more systems or components to exchange information and to use the information that has been exchanged" [24].
With reference to the platforms for smart and physical space related applications (as also telemedicine applications), based on [25], the clear distinction between the following separated conceptual interoperability levels should be considered: communication level, service level, and information level interoperability.
Specifically, information (or semantic) interoperability is the shared understanding of information meaning. Service level interoperability is the ability for a system to share, discover and compose services. Communication level interoperability is interoperability at OSI levels 1, 2, 3, and 4 [26].
As already mentioned, this work focuses on information level interoperability and investigates the adoption in healthcare of an interoperability platform which is agnostic with respect to the service and communication levels used. The main benefit of the proposed solution is the possibility to easily create an open and dynamic smart environment where different actors and systems cooperate on the same information store to enrich the shared knowledge.
Scenario for Person-Centric Health Management
The envisioned scenarios of high-quality and sustainable person-centric health management are clearly multidomain in nature, as they address the patient, the doctor, and the scientific community. Information level interoperability is an important concern since row data originates from heterogeneous devices that are inherently incompatible due to lack of standardization and because they are produced by a competing industry. Our vision is also shared by the partners of CHIRON, a 2010-12 EU project currently in progress. CHIRON's goal is to design an architecture solution for effective and personcentric health management along the complete health care cycle. The main challenge is to integrate the most recent patient information and their historical data into personal health systems and to transform collected information into valuable support for decision making [7]. To this end, several requirements should be met. First of all, data need to be gathered from multiple heterogeneous sources, (e.g. sensors, archives, and databases). Furthermore, in order to make them available in a meaningful and easy way, such data should be stored and uniformly represented in a data sharing platform. Once data are stored, some mechanisms to retrieve and analyze such data are needed. Finally, a shared metadata for information representation is required to harmonize the process of assessing the patient clinical situation and to support the doctor in his/her decision making process.
In order to meet these requirements, CHIRON reference architecture defines the following three layers: the user plane, the medical plane, and the statistical plane.
(i) The user plane concerns interactions by and with the patients (monitoring and local feedback). (ii) The medical plane concerns interactions by and with the doctors (assessment of clinical data, diagnosis, treatment planning and execution, and feedback to the patient). (iii) The statistical plane concerns interactions with medical researchers (external knowledge management).
The integrated care cycle is based on the continuous interaction and exchange of information between the abovedefined three planes. An innovative interoperability platform developed within the framework of SOFIA, another ARTEMIS JU project, may implement the core of CHIRON middleware layer. The proposed platform will contribute to the integrated care cycle, by managing an information space of portable and stationary sensors data, supporting different and heterogeneous measurement devices and actuators, and enabling personal profiling and a personalized approach in healthcare.
Information Interoperability and Interoperability Platform
This section describes the platform that ensures interoperability to the addressed health management scenario. The expected benefit is information integration to increase the knowledge about the patient and to facilitate the exchange of information between user, medical, and statistical planes.
The proposed platform consists of a set of tools supporting application development and application execution. It is open source [27] and-as a tool chain-it is called open innovation platform (OIP), while its runtime interoperability support is called Smart-M3 [28,29]. Smart-M3 is based on the following concepts: information consumers and producers are decoupled and all the relevant information is stored on a shared information domain accessible by all the main actors. This approach opens an easy way to new unanticipated cross-domain applications as consumers do not need to know the details of the producer, but they only need to know how to access the shared information base which is based on an interoperable data model and shared semantics.
The OIP intends to introduce in this way a radical change to the traditional application scenario based on fixed business boundaries, as it inherently supports applications that can interoperate independently from their business/vendor/manufacturer origin.
The shared information domain accessible and understood by all authorized applications is called the smart space (SS). SS information is about entities existing in the physical environment, that is, the users, the objects surrounding them, the devices they use, or about the environment itself. Smart-M3 and the entire platform are developed with 15 principles in mind [30]. Out of these, the most relevant for our discussion are the agnosticism, the notification, the evolvability, the extensibility, and the legacy principles. The agnosticism principle states that the OIP should be agnostic with respect to use-cases, applications, ontology used to represent the information, programming languages, services and communication layers exposing the SS and hosting system. The notification principle ensures that clients may subscribe to be alerted upon specified events. The evolvability principle envisions that the OIP should provide means to support clients that adapt (i.e., "dynamically reconfigure" when changes in the SS information space occur). For example, if new relevant sensors are added to the SS, the application should benefit without the need to change the existing code.
The extensibility principle calls for the inherent and efficient support to add features to the OIP, both at information and service level. Based on wise tradeoffs between the extensibility and the agnosticism principle, appropriate ontologies and data models could be defined, for example, to add access control and privacy policies and enforce authentication. Information security and privacy, as well as trust management, are fundamental qualities in a telemedicine scenario, that may become an OIP extension. The legacy principle states that existing (i.e., legacy) devices may exchange information with the SS with interface modules called "legacy adapters." Smart space and legacy adapters and their architectural support in the OIP are conceptualized as follows.
The OIP is built around its interoperable data sharing framework, named Smart-M3, and released by Nokia within SOFIA project. Smart-M3 defines two types of entities ( Figure 1): the semantic information broker (SIB) and the knowledge processor (KP).
The SIB is a digital entity where relevant real-world information is stored and kept up-to-date. It is an information store for sharing and governing all the data relevant for the domain of interest. The information model is a directed labeled graph corresponding to a set of RDF triples (resource description framework a basic semantic web standard [31]). The information semantics is specified by ontologies defined in OWL [31]. A query language [32] augmented by an inferential component provides reasoning support to applications.
The KPs are software components interacting with the SIB and producing and/or consuming data. The legacy adapters are KPs that enable legacy SS-unaware devices to exchange information with the SIB (Figure 1). A KP exchanges data through a simple protocol named smart space access protocol (SSAP), an application layer protocol based on XML. The SSAP defines a simple set of messages (join, insert, remove, update, query, subscribe, and leave) that can be used over multiple connectivity mechanisms, such as TCP/IP, Bluetooth, OSGI, and NoTA [25]. The join is the first operation done by a KP in order to register itself to the smart space. After this operation, the KP can write or delete a subgraph (i.e., a list of triple) using the insert or the remove primitives. The update is an atomic operation including an insert and a remove message. The KP can retrieve information from the SIB with the query operation. The SIB supports a number of different query languages including triple-based queries and Wilbur queries [32]. The subscribe operation allows a KP to specify conditions at information level that, when verified, are notified automatically according to the subscription notification paradigm. When a KP subscribes to a part of the graph, it receives a notification whenever such graph is modified; since the notification message is provided by the SIB with a short delay and contains details about the modifications that occurred, the KP is able to react promptly and in a specific way. The leave operation is used by a KP to specify the end of its session with the SIB. A KP that performs the leave has to join again in order to interact again with the SS. KPs may be developed using an application development kit (ADK), which is meant to increase developers' productivity by hiding the ontology and the SSAP protocol details, raising in this way the programming level of abstraction. Applications may be developed in several popular programming languages, including Python, ANSI C, C#, and Java.
Scenario Description
The addressed scenario consists of users living and moving in an environment monitored by environmental sensors. The user also relies on personal devices, among which are medical devices, both wearable and not wearable. The user moves freely inside this smart environment and all health and user's context data are continuously monitored and collected.
The first requirement to build a smart space out of this scenario is to create a shared and comprehensible digital description of all the entities that play a significant role in such an environment. In order to map each physical entity to its digital representation in the smart space, each physical entity needs to be uniquely identified. After that, each physical entity and its properties should be described with an appropriate data model. This description is created at smart space initialization time and concerns the users (both patients and doctors), their profiles, the devices, the environments, and all the relations between these main actors. Relations between entities may change dynamically and this may be reflected in the smart space manually using GUIs or automatically using some identification or localization technology. For example, in the proposed clinical scenario, the association of the same set of devices to different patients can be done explicitly with a dedicated GUI or it may be done automatically using the RFID technology as follows. The medical devices and an RFID reader are located on a desk. The patient registers himself just putting his RFID tag on the RFID reader; then all measures taken by the devices on the desk are associated to this patient until the RFID tag is removed.
The architecture of the proposed system is based on the platform described in Section 3 and it is shown in Figure 2. The OIP supports the user plane and its access from the medical plane. Different software agents (KPs) cooperate through the semantic information broker.
In the user plane, data from heterogeneous sensors are collected by a PC or a smartphone and sent to the SIB. The sensors adapters are legacy adapter KPs (Figure 1) that could be seen as ontology driven translators: they enable the exchange of information between the SS and the legacy world.
The SIB is the core of the system; it stores and shares not only the data received from the devices but also all the information created during the initialization process. This implementation is consistent with the user profile concept proposed in CHIRON and called the "Alter Ego" [7]: the "Alter Ego" is an evolving virtual entity modeling all the relevant aspects related to the user health including medical history, habits, and functional health status. This entity has the capability of evolving and adapting over time to various domains and user conditions. Once generated, the profile is kept up-to-date automatically through the information provided by the multisensorial platform; in any case the doctor can add his inputs at any time. The user profile includes static/semistatic and dynamic parameters.
Based on the data collected by the sensors and on the knowledge of the relations between all the entities involved, new services can be devised. With reference to Figure 2, some examples of these services follow.
Aggregators are generic services consuming information from the SIB and enriching such information according to specific inference rules. Enriched information is stored back in the SIB in order to increase the knowledge and it includes, for example, indexes, aggregated parameters, or new semantic connections between existing entities.
The automatic alarm generator is a service that generates an alarm condition on the patient profile as soon as at least one patient parameter goes out of the range prescribed by his/her profile. The conditions for raising an alarm are set by the doctors for every patient and are built into the patient profile.
History service collects relevant data so that it can be statistically analyzed to discover or validate macroscopic relationships between patient profiles, diagnosis and the efficacy of the medical treatments.
The alarm dispatcher service notifies an alarm status to the appropriate doctor using standard communication services like SMS and e-mail.
With reference to the medical plane, the following are examples of policies that could be implemented: medical staff could visualize all data collected from his/her patient through different platforms (PC or smartphone) in real time. Doctors could monitor their patients and specialists could be notified in case of alarm. Doctors could visualize, modify, and set a manual alarm according to a rule based-policy.
Starting from this scenario an application linking the user plane with the medical plane was implemented.
Application Design and Implementation
The implementation of a smart space application based on the proposed approach first requires an ontology to describe the domain of interest, then a set of KPs need to be developed in order to achieve the desired behavior.
An ontology is a formal representation of the domain of interest made up of concepts and relations between them. When approaching ontology design, the designer must define classes, arrange them in a taxonomic (subclasssuperclass) hierarchy, and define properties with their features (e.g., functional, inverse-functional, and symmetric). There is not a unique and optimized way to model a domain. On the contrary, there are always viable alternatives and the best solution depends on the application requirements and should be as intuitive, extensible, and maintainable as possible [33]. The ontology used in this application is an extension of an ontology modeled for a previous demonstration [34]. The ontology class tree is shown in Figure 3. The main classes are Person, Environment, Data, Device, and Alarm; they are all subclasses of Thing.
The Person and Environment entities are self-explanatory. Devices are entities that can produce data or run KPs and are described by their characteristics, (e.g., MAC address, protocol). A Data entity is described by a value, a Measurand, a Unit Of Measure, and a timestamp. By modelling the data class in this way, we ensure that any KP that consumes sensor data would be able to take advantage of new sensor types without having to rethink the KP. Alarms are entities characterized by an AlarmType, for example, HeartRateAlarm.
In our application, we have some users moving in an indoor space divided into two rooms and each environment is monitored by the sensors of a wireless sensor network (WSN); a ZigBee WSN, developed by Eurotech, is used to sense temperature, humidity, and the presence of water on the floor. In order to demonstrate the multivendor interoperability of the system, Intel iMote2 nodes were added to each room to further sense their environmental conditions. Both Eurotech WSN and iMote2 nodes send data to a local PC (i.e., the user's home PC) where there are two KPs, feeding the SIB, that is, one for the Eurotech WSN and one for the iMote2. The amount of sensors and the sensor network configuration do not need to be known a priori: in fact the system is able to benefit from certain sensors even if added at runtime. According to our ontology, when a sensor is added to the system, its KP inserts in the SIB a new instance of Data semantically related to the monitored entity; in this way, consumer KPs can discover all the data associated to a certain entity with just a query.
In our user case, we plan to do out-of-hospital monitoring of patients with a cardiovascular disease. The user is provided with the following medical devices: an A&D scale (UC-321PBT), an A&D sphygmomanometer (UA-767PBT), a Nonin finger pulse oximeter (Onyx II 9560), and a Zephyr BioHarness, which is a smart wearable device capable of collecting several vital signals.
All of these devices communicate via Bluetooth, each with its specific protocol. To satisfy the usability requirement and simplify user interaction, the user only needs to turn on the devices and take the measurements; data are sent to the repository without any further action.
We achieve this functionality by exploiting the initial knowledge about the environment shared in the smart space. The only information needed, in fact, are the device's communication characteristics (MAC address and protocol) and a semantic connection associating the device to its current user. There is only one KP associated to each patient. This is the PhysiologicalSensors KP which communicates with all nearby Bluetooth devices semantically connected to the patient, and sends their data directly to the SIB. This KP works as follows: once launched, it queries the SIB for all information about the devices related to its user including their MAC addresses; then when one of these devices is discovered, it implements the proper protocol, parses the data, relates the data to its patient and updates the SIB. In the current implementation, this KP runs on a smartphone so that the sensors are wearable and the user can move while taking the measurements. In a different way, as shown in the previous section, a sensor can be stationary, shared, and only temporarily associated to a patient, for example, through the patient's RFID tag. While the RFID reader KP is responsible for handling the dynamic "semantic connections" between patient and instruments, the patient's PhysiologicalSensors KP dynamically searches all the associated sensors and gets their data as described.
As user location needs to be known, each room is equipped with an RFID reader located next to the room entrance, while, as previously mentioned, each user has an own RFID tag for example attached to his/her smartphone. When a person enters a room, the corresponding tag is read and at information level the user is associated with the proper environment by the Location KP. All information sent to the SIB can be used by multiple applications. To demonstrate the proposed personcentric health management approach, the following agents were implemented: the ThomIndex, the AlarmGenerator, the AlarmNotifier, the HistoryLog, the HealthCareMonitoring and the ManualAlarmGenerator.
The ThomIndex KP evaluates the Thom Index (TI) which is a bioclimatic parameter measuring the perceived discomfort level originated by the combined effect of the environmental humidity and temperature [35]. The TI is evaluated as follows: the SIB is searched for all humidity and temperature data available; then the TI for each room is calculated based on the mean temperature and humidity values provided by the available and properly working sensors; eventually the TI associated to each room is stored back into the SIB. This KP is an example of "Aggregator KP" as it has the ability to aggregate raw data into higher level data with an increased semantic value.
The AutomaticAlarmGenerator KP is meant to publish the SIB alarm conditions detected according to properly specified policies. Currently, alarm conditions are very simple as they are just threshold based. If a user has a safety threshold inserted by the doctors and associated to certain parameters, then the KP performs appropriate subscriptions to the smart space in order to be notified of all updates of the sensible parameters. In Figure 4, a subset of the semantic graph used by the AutomaticAlarmGenerator KP is shown. Using the notation adopted in [36]: (i) the classes are represented in orange; (ii) each instance is connected to the classes it belongs to by the red dashed arrows representing the rdf:type property; (iii) classes, instances, and properties connecting them are uniquely identified by URIs (universal resource identifiers) whose semantic is commonly understood by all the software agents; (iv) literals (i.e., numerical and text values of properties noncorresponding to URIs) are shown in green. KPs access the graph by navigating through it with query operations and use the semantics of URIs to interpret the meaning of data. In Figure 4, the Person Irene has a safety threshold on her heart rate with maximum value 150 and minimum 50; as her data measuring heart rate exceeds the maximum value, an instance of alarm is generated and KPs subscribed to this information may find all the relevant data to react promptly to this situation. In our ontology, in fact, the instances of the Person class are also related to their respective doctor which can be contacted by another important KP: the AlarmDispatcher.
The AlarmDispatcher KP subscribes to the creation of Alarm instances, and, as soon as it is notified, it sends an email or an SMS with all related information to the appropriate doctor. Doctors' profiles include relevant personal information, particularly their email address and phone contact. Two alternative SMS sending solutions were implemented: the first runs on a Nokia N900 smartphone, while the second one on a PC connected to an embedded Siemens HC25 radio module. The History KP subscribes to all user data variations and logs the notified information in a file with an associated timestamp. The above KPs demonstrate our user plane implementation.
Moving to the medical plane, this is currently handled by a single KP named HealthCareMonitoring KP, which tracks all users' properties and allows the health care service to monitor a patient in realtime.
The KP allows the healthcare staff to select a patient, then creates a subscription to all of his/her relevant data. Therefore, it shows in real time all his/her data stored in the SIB, so not just physiological data, but also the person's environmental information as well, that is, health parameters together with the Thom Index and the environmental conditions of the place where she is located.
So far, based on the available sensors, the following health parameters may be collected: heart rate, respiration frequency, skin temperature, activity index, posture angle, weight, oxygen saturation, diastolic, systolic, and mean pressure. The HealthCareMonitoring KP detects alarm conditions and alerts the medical plane on top of the user plane AutomaticAlarmGenerator KP. Figure 5 shows the GUI of this KP. Furthermore, the healthcare staff, who monitors the patient's profile and current data, can also generate a manual alarm using the ManualAlarmGenerator KP. The interaction between the implemented KPs is depicted in Figure 6 while Figure 7 shows some preliminary data collected in our laboratory from a user wearing his BioHarness and walking around two rooms. The graph in Figure 7 shows the history of the following parameters over a short period of time during a hot summer day: environment Thom Index and patient's heart rate, respiration rate, and skin temperature.
Conclusions and Future Work
In this paper, a framework to support person-centric health management was presented and the proposed approach was demonstrated with a use case that focused on real-time monitoring of patients' health and environment. In order to handle the heterogeneity of relevant information and to overcome the fragmentation of the instrumentation involved, a shared and open source interoperability component was proposed, which is ontology driven and based on the semantic web data model. A significant set of multivendor physiological and environmental sensors was considered, and their data were collected, processed, and monitored.
With the proposed approach, information consumers and producers are decoupled and relevant information is stored on the shared information search domain offered by the interoperability platform and called "smart space." The application functionalities are based on the cooperation of different software agents exchanging information through the smart space. Agents run on multivendor platforms based on different operating systems and written in various programming languages.
Smart space interoperability makes clear separation between device, service, and information level interoperability. The described work was focusing on information level interoperability, so primary requirements such as privacy and security were disregarded, being considered as service level qualities, therefore not relevant to the purpose of this paper. Consequently, the proposed solution now needs to be integrated within an SOA (i.e., a service-oriented architecture) and the plan is to integrate it within the architecture being developed within the already-mentioned CHIRON project. Here, standard ontologies will be adopted and the entire patient history together with the external information provided by the statistical plane will be offered in a unified view as a contribution to win the healthcare management challenge. | 7,897.4 | 2011-07-24T00:00:00.000 | [
"Computer Science",
"Engineering",
"Medicine"
] |
bond market analysis: tHe main Constraints in tHe researCH of 21 st Century
. Searching for alternative source of bank financing, the view on capital market is taken. Recent research on capital market issues are arranged into four di -mensions: theory and assumptions of efficient capital market, government’s role in it, other distortions and global interrelatedness. Main investigations are decentral ized and visualized in “theoretical eight” model. Conclusions made on the diver-sity of interpretation of market efficiency, strongly expressed demand of informa tion symmetry, soft actions of governments and the value of foreign performance in domestic markets. Furthermore, new approach to the classification of countries by their maturity in capital market is argued. The state of art of 2009-2012 of bond market and government debt is briefly described.
Introduction
Making a brief surveillance of financial market, capital market and its sub-marketbond market -are identified as key areas, arguing by the most attractive risk free or least risky investments during economic shocks, serving for financial resources instead of banking, growing demand for the government debt.
sought to quantify the impacts on capital and especially bond markets, but the literature is still relatively sparse.
Therefore the purpose of this paper is to summarize issues on which 21 st century research on capital and its sub-sector -bond market -is based and markup the direction of further investigations needed.
The paper contains of five sections. The first section gives smooth introduction into a paper. Section No 2 is briefly summarizing most research dimensions, such as theory and assumptions of efficient capital market, government's role in it, other distortions and global interrelatedness. Section No 3 provides new approach to the classification of countries by their maturity in capital market. New classification is argued and visualized. Section 4 contributes to description of state of art of current capital market, country gross debt situation and their discrepancies. Conclusions (Section No 5) are made.
21 th century capital market research evolution
Mostly analyzing the scientific literature on capital market issues the effective market theory is being introduced and investigated. Assumptions of perfect capital market are being analyzed by Modigliani and Miller (1958). All financial claims are perfectly divisible, there are no transaction costs, there are no taxes, market is competitive and other conditions form the perfect capital market (Ho, Bin Lee 2004). However, marked conditions are of rare existence in financial markets.
Therefore, modern theory of finance is introducing different analytical approach to effective capital market (Leipus, Norvaiša 2003). As the main function of capital market is the distribution of assets and equity, the effective role is being formed by price model. The description of effective capital market hypothesis emphasizes the correction of prices, which reveal the right information. Therefore the effectiveness of capital market is being explained as the symmetric and right information, possessed by market players (Fama 1970). To sum up, the effective capital market model was evolving from couple perfect market conditions to information symmetry as the main constraint of price forming.
Another approach to capital market's effectiveness is macroeconomic. It is measured through its effective capital distribution (Pekarskienė, Pridotkienė 2010) by enabling growth of GDP. The authors structure the research on the conflicts and stresses of efficient market hypothesis in market of securities. The conclusions highlight the argument that the national securities market is much more influenced by the processes of globalization than just country-specific economic developments. Although the impact of globalization on all spheres of life is certain, however, there is dependence between the stock market activity and the economic situation of the country. s tudies have shown that nearly in the majority of countries there are relationship between stock price indices or other characterizing capital market indicators and gross domestic product, inflation rates (Pekarskienė, Pridotkienė 2010).
However the contradiction presented denies or diminishes the role of current capital market effectiveness by emphasizing correlations between different financial and macroeconomic variables. Continuing the analysis of objections, Stankevičienė and Gembickaja make the notion on efficient capital markets by the behavior of a rational investor. Over the past few years, from the investor's point of view, the vulnerability of the markets has led to increased uncertainty and unpredictability, as market conditions cannot always be judged with the help of standard financial measures and tools. Despite strong evidence that the stock market is highly efficient, i.e. one cannot earn abnormal profits by trading on publicly available information, there have been a number of studies documenting long-term historical anomalies in the stock market that seem to contradict the efficient market hypothesis. During the recent years, the examples of market inefficiency in the form of anomalies and the irrational behavior of the investor have been observed more frequently. The existing phenomenon can in part be attributed to the less-than-rational aspects of investor behavior and human judgment (Stankevičienė, Gembickaja 2012). To sum up, the capital market effectiveness could be interpreted through several approaches: classical hypothesis of effective markets, information that forms price model, macroeconomic access as well as behavior of its participants (e.g. investors).
Last decade's investigations on the topic of capital market could be divided into 4 dimensions: constructing and analyzing the efficient capital market model or hypothesis, evaluating the impact of such distortions as information asymmetries on capital market, comparing interconnectedness of global capital markets, measuring the government's role in the capital market, arguing on its regulation and guarantee giving's, shown in Figure 1.
Figure 1 describes theoretical "eight" -dispersion of research topics. The axes show the dimensions into which all recent research works by international organs (IMF, WB, OECD, different institutes and universities) are divided. The more distant the ball with names of authors is the higher intention on the related topic it introduces. Full balls represent the dependence on both topics in vertical and horizontal axes, others belongs to one domain where the position it's taken.
All dimensions could be combined into life-cycle of efficient capital market (following the Fig. 1 anticlockwise), which could be reached by diminishing the impact of such obstacles as information asymmetries, including government role with acceptable restrictions and getting interconnected in global market.
In the relation with the first paragraph, other efficient capital market models and hypothesis are being analyzed in combination with opposite approaches such market distortions as information asymmetries. Klimašauskienė and Mosčinskienė analyzing the effectiveness of capital market state that information efficiency in country specific capital market occurs on its light form (Klimašauskienė, Mosčinskienė 1998). As well as Leipus and Norvaiša, Klimašauskienė and Mosčinskienė emphasized the significance of securities price. The findings were based on information effectiveness which takes its evidence when security buyers and sellers share the same data and similar expectations reflected to market price. Rudnicki agrees that there occur some events that contradict the efficient market hypothesis therefore they are called anomalies (Rudnicki 2012). By analyzing stock splits, the author indicates their implications on narrowing the information asymmetry between managers and shareholders as well as diminished probability of informed trading.
Introduction of the derivatives as financial innovations into capital market quantitative analysis showed that the effectiveness of macro-prudential policy in this environment depends on the government›s information set, the tightness of credit constraints and the pace at which optimism surges in the early stages of financial innovation. The policy is least effective when the government is as uninformed as private agents, credit constraints are tight, and optimism builds quickly (Bianchi et al. 2012). These arguments state the opinion of transition of effectiveness through related markets as well as their constraints. While being well informed government forms prudent policy and macroeconomic indicators rise, the capital market is functioning effectively. Conclusion is being approved by Yuko Hashimoto and Konstantin M. Wacker, who investigate whether better information about the macroeconomic environment of an economy has a positive impact on its capital inflows, namely portfolio and foreign direct investment (Hashimoto, Wacker 2012). Moreover the informational imperfections in credit market are described as 'micro" concern relating consequences. In this line of inquiry, problems of asymmetric information between borrowers and lenders lead to a gap between the cost of external and internal financing. This notion of costly external financing stands in contrast to the more complete-markets approach underlying conventional models of investment emphasizing expected future profitability and the user cost of capital as key determinants of investment (Hubbard 1998). While watching positive correlations between macroeconomic indicators and information symmetry (analyzed by b ianchi et al. 2012;Hashimoto, Wacker 2012), the main causes of imperfection are found in microeconomics (Hubbard 1998).
The information asymmetry is analyzed in venture capital dimension as well. Krystyna Brzozowska initiates that establishment of separate funds in each region's appropriate environment where the rules of venture capital investments are well known. On the other hand, venture capital funds will have difficulties in monitoring their investee companies as well as to provide suitable advisory capacity to them. It can be assumed that inexperienced management team in most early stage of ventures and the greater uncertainty (less information) connected to new technology and market will evolve some problems difficult to solve (Brzozowska 2008). Meanwhile searching for causes of information asymmetries in microeconomic factors or even behavior of market participants, here by the new factor of technology is introduced. As the market progression is formed by technology and information share is facilitated, this approach is being questioned and requires more arguments and research.
To sum up, the results of the research on capital market effectiveness and information asymmetry dimensions could be summarized as capital pricing model: − Effective capital market hypothesis emphasizes the correction of prices, and is explained as the symmetric and right information, possessed by market players.
− Symmetric information correctly allocates the resources, regarding the main capital market function of distribution.
These could be described as highly correlated issues of the capital market theory: reaching the perfection while eliminating the imperfection of information shared in the market.
Another dimension which is being discussed, the role of government, its regulations and guarantees in the capital market. Concerning government guarantees on bank bonds they were adopted in 2008 in many advanced economies to support the banking systems. They were broadly effective in resuming bank funding and preventing a credit crunch. The guarantees, however, also caused distortions in the cost of bank borrowing (Grande et al. 2011). Contrarily to guaranties, governmental restrictions are being valued. Claudio Raddatz and s ergio L. s chmukler studied the relation between institutional investors and capital market development by analyzing asset-level portfolio allocations of Chilean pension funds between 1995 and 2005. In the analysis of government restrictions, the conclusion was come to pension funds may have contributed to the development of certain primary markets, but not a force driving the overall development of capital markets, because of asset illiquidity and manger incentives (not regulatory restrictions) (Raddatz, Schmukler 2008). Authors indicate the absence of impact of governmental restrictions on overall market development while influence on primary markets (e.g. capital) is being identified. Therefore the conclusion has come of action based rebound on action targeted market. Though Meng and Pfau find a significant impact of pension funds on capital market development in the overall sample, this result is driven by countries with "high" financial development (e.g. United States, United Kingdom, Japan, Germany, etc). For countries with "low" financial development (Argentina, Peru, Poland, Hungary, etc.), pension funds do not show a significant impact. Countries with different levels of financial development have different financial market climates that can directly impact the role and performance of pension funds. Differences include pension fund investment regulations, market efficiency, transparency, the legal framework, market activities, and macroeconomic and financial conditions (Meng, Pfau 2010).
Several working papers are analyzing the government's role and interconnection issues in combined way of research. N. E. Magud, C. M. Reinhart, K. S. Rogoff have constructed two indices of capital controls: Capital Controls Effectiveness Index (CCE Index), and Weighted Capital Control Effectiveness Index (WCCE Index). It was found that there should exist country-specific characteristics for capital controls to be effective. Capital controls on inflows seem to make monetary policy more independent, alter the composition of capital flows, and reduce real exchange rate pressures (although the evidence there is more controversial) (Magud et al. 2011). It follows that governmental regulations re-charge the structure of capital flows without a movement trend.
Meanwhile P. S. Srinivas, E. Whitehouse and J. Yermo compare the rules in the new systems of Latin America and Eastern Europe with richer OECD countries on regulating the pension fund industry's structure. It was revealed that taxonomy of investment risks in pension funds is a light limit on domestic investments (on equities and bonds) (s rinivas et al. 2000). It comes to conclusion of regulations limiting domestic capital flows. Other expression of government actions is represented by L. Jaramillo and A. Weber. They have come to conclusion that fiscal variables do not seem to be a significant determinant of domestic bond yields in emerging economies. However, when market participants are on edge, they pay greater attention to country-specific fiscal fundamentals, revealing greater alertness about default risk (Jaramillo, Weber 2012). The interpretation is made about the governmental role which is assessed more ponderable to domestic financial instruments than on capital market development in the overall sample. The regulations seem to be absent of significance in capital market.
On the other hand S. J. Peiris assume foreign participation in the domestic government bond market. Author states that greater foreign participation tends to significantly reduce long-term government yields. Moreover, greater foreign participation does not necessarily result in increased volatility in bond yields in emerging markets and, in fact, could even dampen volatility in some instances. However, foreign investors could act as catalysts for the development of local bond markets, particularly by diversifying the institutional investor base and creating greater demand for local emerging markets debt securities. The author comes to conclusion that institutional investors, both domestic and foreign, have played a critical role in developing capital markets in most mature markets and in more developed emerging markets (Peiris 2010). The division of country profile in capital market could be mentioned: emerging and more advanced economies are being analyzed.
Moreover, bond yields and other securities' prices are a common goal or even an issue for investigators of governmental regulations impact as well as those who analysis other market distortions and spillovers in capital market. Alex Sienaert finds more disadvantages of foreign invasion by examining the causes, nature and impact of rising participation of foreign investor in local currency bond markets of developing country. Much of the volatility in returns occurred through the currency channel (not bond prices in local currency), insulating dollar-hedged and local currency-benchmarked investors. An important source of selling pressure on emerging markets local currency bond markets was forced by liquidations of foreign investors due to the relatively low collateral value of emerging markets bonds. The growing interconnectedness of global capital markets increases the sensitivity of emerging markets asset prices, including bonds, to global factors (s ienaert 2012).
Furthermore, supervisors and regulators cannot develop the markets directly; only borrowers and lenders can do this. This distinction is not always appreciated, and governments at times go too far in their efforts to facilitate financial market development. This is apparent in the most common strategy for government-led financial market development, which is the "Build It and They Will Come" approach. In this approach, the government introduces not only the legal infrastructure but also particular instruments and exchange mechanisms, in the expectation that private players will rush into the ready-made markets. The problem in many cases is that few agents actually come to play and often there is limited activity in these new markets (Chami et al. 2009).
Further investigations are made on global interventions of capital markets and the local influences of it. Tamim Bayoumi's and Trung Bui's identification through hetroscedacity to estimate spillovers across U.S., Euro area, Japanese, and UK government bond and equity markets in a vector autoregression. The results suggest that U.S. financial shocks reverberate around the world much more strongly than shocks from other regions, including the Euro area, while inward spillovers to the U.S. from elsewhere are minimal. There is also evidence of two-way spillovers between the UK and Euro area financial markets and spillovers from Europe to Japan (Bayoumi, Bui 2012). Other researchers analyze financial integration or interrelatedness, asking how Asia compares with Europe and Latin America (Eichengreen, Luengnaruemitchai 2006). Spillovers all across the world are identified and the impact to local capital markets is agreed. Domestic and external factors on performance of capital market are analyzed by Irina Bunda, A. Javier Hamann, and Subir Lall. The co-movement in emerging market bond returns and disentangles are of influence of external and domestic factors. The conceptual framework, set in the context of asset allocation, allows describe the channels through which shocks originating in a particular emerging or mature market are transmitted across countries and markets (b unda et al. 2010). Common shocks or to "pure" cross-country contagion of spillovers is accepted.
Summarizing the consequences of governmental actions (guarantees and regulations) and global interconnectedness separately and in together, the conclusions could be drawn: − An advantage of government restriction is the development of certain primary markets, but not a force driving the overall development of capital markets.
− The drawback of governmental regulations is re-charged capital structure without a movement trend, limiting domestic capital flows.
− Foreign presence in local capital market results proses (increased volatility, catalysts for the development, greater demand) and coins (increased sensitivity of asset prices) depending on maturity of capital market.
Developing capital market: 4 stages
According to issues (market distortions, effective market, government regulation) and regions (OECD countries, Latin America, Europe, Asia) examined in section 2 the stages of capital market development could be distinguished. All researches agreeing on imperfection and the sequence of development of capital market, countries are divided into emerging and matured.
The authors of this paper propose the following division into 4 groups with the arguments below: First argumentation comes from the agreement on opinion that country's economic conditions are determined by the behavior of financial markets. For example, it is said, when prices of shares are starting to fall, one can expect the economic stagnation. And vice versa, the growing trend of share price alerts the potential economic growth (Leipus, Norvaiša 2003). The health of financial intermediaries and markets is crucially dependent on the health of the private and public sectors. There are potentially many types of infrastructure that governments need to build: the legal system, including bankruptcy procedures; a modern payment system for clearing and settling securities transactions, retail payments, and large-value payments; instruments, in the sense of legal definition and recognition; and markets, including rules and possibly physical infrastructure for the operation of primary and secondary markets (Chami et al. 2009). Questioning whether the financial markets directly impact the economy, very little evidence corroborates this view. There are much more signs pointing to the fact that financial markets simply reflect firms expectations on the behavior of the economy in the near future. These "mirrors" is generally regarded as different countries' financial market indices: Dow Jones Industrial Average (DJIA), Standard and Poor's (S&P), OMX, etc (Leipus, Norvaiša 2003).
Most financial institutions (IMF, WB, UNDP) casual inspection suggests that currently the classification systems which are quite similar in terms of designating countries' economies as being either 'developed' or 'developing'. Given the large and diverse group of developing countries, all three organizations have found it useful to identify subgroups among developing countries (Nielsen 2011). Therefore remaining matured capital market range, the emerging one could be divided into three subsections: underdevelop, emerging, integrated.
In the global capital market statistics still there are countries lacking attention on global research and data, which capital market share (% of GDP) is too miserable to be watched. There are many markets in which borrowers and lenders are present, the instruments used are agreeable to both parties, and yet the market has little activity beyond primary issuance and redemption. For example, many of the nascent government bond markets around the world are simple "buy and hold" markets. While such markets help achieve the fiscal policy goals of the government, they do not lead to financial market deepening and its accompanying benefits. This is because there is no trading in the instruments, and in particular, no agents making a secondary market in the securities (Chami et al. 2009). In some EU countries, especially new members, capital market has not yet been developed, but it is still growing together with the increase of the level of innovation and entrepreneurial activities (Brzozowska 2008). Therefore a level before the 'emerging' classification should be drawn in order to provide comprehensive analysis and ensure the sustainable development of capital market, which could be reasonable as fine segmentation specializes in targeting measures to successful achievement of the goal.
According to sequencing approach and logic, the country should have intermediate period from performing in emerging market rating to becoming a matured one. Financial market development is seen as both the wider use of existing financial instruments and the process of creating and adopting new financial contracts for intermediating funds and managing risk. A key aspect is that development occurs when market players are able to reach mutually acceptable compromises regarding the terms of financial transactions. Agents strike grand compromises, such as those between maturity and collateral, and between seniority and control, as well as myriad smaller ones (Chami et al. 2009).
Criteria for reaching, valuing the acceptance of classification level could be defined as follows: the requirements for capital market development, and compare and contrast experiences across both mature and emerging, under-developed and integrated markets -benchmarking, corporate governance and disclosure, credit risk pricing, the availability of reliable trading systems, and the development of hedging instruments. These are fundamental for improving the breadth and depth of corporate debt markets (Luengnaruemitchai, Ong 2005). Moreover, as it was mentioned before, regulation and supervision play a supplementary role in market development. An important job of the regulator is to establish a supportive infrastructure for contract enforcement and dispute resolution. This infrastructure has many concrete as well as abstract features, but collectively these aspects have come to be known as the "rule of law" (Chami et al. 2009). Furthermore, the presence of derivatives can quicken market development in the underlying, and if the infrastructure and regulatory framework is available, their introduction need not be delayed. The policy challenge is to support the creation of an intersection between the set of desired instruments and the set of feasible instruments, and to enlarge it over time. Often, this intersection must be created by eliminating or overcoming obstacles that prevent an instrument from being introduced or used (Chami et al. 2009). Briefly, the view of financial market dynamics can be expressed as follows: if borrowers and lenders are willing and able to contract, and liquidity providers find conditions conducive to trading the instruments that are created, then financial markets will develop. The regulatory structure can support this process by removing obstacles that make potential borrowers, lenders, and liquidity providers unwilling or unable to play their roles, and by creating the right incentives for each agent to fulfill their end of the bargain (Chami et al. 2009).
Forming the sequences of the countries by their capital market development stages more intuitively, the range scale is being of the important need to be implemented for overall, coherent and integral research and data collection. The scale should contain the key variables of capital market (e.g. derivatives) and all market players (e.g. government). Rages should be specified and unified by each scale grade.
Bond Market Development: state of art
Bond market is taken into consideration as the development of a "risk-free" asset is a key step in financial market development. The government is often thought of as the entity with the lowest credit risk in an economy (Chami et al. 2009). However, for example, in Hungary, one of the earliest corporate borrowers was the local subsidiary of McDonald's Corporation, which was widely perceived to have a better credit rating than the government.
Since the mid-1990s, corporate bond markets have become an increasingly important source of financing for the private sector, especially in the emerging market countries. The authorities in these countries are becoming increasingly aware of the importance of establishing deep, liquid corporate debt markets and have placed such development high on their agenda. To date, corporate bond markets in many countries remain largely underdeveloped, with a limited supply of quality issues and inadequate market infrastructure. Even in mature market countries, such as the United States and Europe, secondary markets for corporate bonds are relatively illiquid for the majority of bond issues, in the same manner that liquidity in government securities markets is usually limited to a few benchmark issues (Schinasi, Smith 1998).
Joshua Felman et al analyzing the Asian bond market development have overcome with reasons for capital and its symbiosis bond market to develop that could be applied more widely (globally): − Finance systems are extremely bank-centric, which meant that most of the financial risks were being concentrated in the banking system -and there is no alternate channel of intermediation that could be used if the banks once again encountered difficulty (Felman et al. 2011). The Baltic States, Central European countries could be specified likewise. Although there is no definitive evidence that either a marketbased or bank-dominated financial system is better. However, it has been argued that a more diversified financial system would mitigate its vulnerability to systemic risk. For instance, the effects of the Asian crisis and the recession in Japan during the 1990s may well have been far more benign if the countries involved had well-functioning capital markets and correspondingly less heavy reliance on their troubled banking sectors during this period (Luengnaruemitchai, Ong 2005). Moreover, the relative unimportance of the corporate bond market in Europe was mirrored by the corresponding dominance of the banking sector. This is in direct contrast to the United States, where banks play a small role in the financing of large companies, and face strong competition from the corporate bond market even for medium-sized companies (Schinasi, Smith 1998). Roldos re-argues that banking and bond markets could be developed in tandem, by building an appropriate regulatory and institutional framework to encompass both. Although local securities markets provide an alternative source of funding to the banking sector, especially during banking crises, a sound and well-regulated banking system could be a necessary and desirable complement to the development of local securities markets (Roldos, Jorge 2004). In emerging markets, it has been noted that the Central European countries have little intermediary capacity to underwrite corporate bonds. The large, foreign-owned banks in these countries have little incentive to devote capital to such activity in the local market, while the local banks and brokerages typically lack the resources to do so. In Thailand, banks have been reluctant to underwrite bond issuances, possibly because they fear competition from the bond market. The opposite is true of banks in Hong Kong SAR, which have begun to underwrite bonds to take advantage of the attractive fees from the process. However, the advent of several banking crises in some of these countries has led to the realization that the sources of corporate borrowing need to be diversified. That said, the corporate debt markets in many emerging market countries remain underdeveloped.
− b orrowing had suffered from a double mismatch, since long-term domestically oriented investment projects were being funded through short-term and foreign currency borrowing (Felman et al. 2011).
− Countries in the region were perceived to be excessively dependent on volatile capital inflows, a situation that struck many observers as ironic since the region had an abundance of domestic saving (Felman et al. 2011).
− The rise of foreign interest in domestic bonds has another important ramification:
growing off-shore activity. Foreign investors are increasingly obtaining exposure to emerging markets by using various "access products", such as over the counter derivatives, structured securities, or offshore special purpose vehicles (Felman et al. 2011).
Describing the volume of nowadays bond market the data of total debt securities of all issuers by country is taken into consideration. As one could see from the Figure 3, countries could be easily divided into classification groups described in section 3. UK, US, Japan, Italy, France and Germany belong to the matured capital markets country group, while Hungary, Argentina, Poland, Russia and others -emerging. There are no representatives for the group of under-developed capital countries as statistic data is too poor to gather. Otherwise Belgium as well as Netherlands is balancing between emerging and matured capital market countries, making it a strong argument to assign them to an integrated capital market group.
The dynamics of the period of year 2009-2012 is slightly: there is no significant gain in development as well as sharp decrease. On the other hand, some structural changes could be foreseen in Japan and USA: during the last four years the borrowing volumes did increase by average of 10 per cent, potentially caused of presence in the top rankings and favorable cost of borrowing. Looking for issues on causes for bond or other debt securities market extension, the gross central government debt data is taken into the consideration (Fig. 4). During the period of the year 2009-2012, the sharpest trends are seen in Japan and United States as well as in Fig. 3. The conclusion of bonds and other debt instruments being the source financing government debt is made. Furthermore, the trends in Figs 3 and 4 dynamics could be remarked. As well as France's , Germany's, Italy's, United Kingdom's debt securities stand out among the other countries in the analysis (Fig. 3), gross government debt in the mentioned countries have the same characteristics (Fig. 4). In later comparison one could note the amount differences between amounts of total debt securities and gross government debt, which could be explained by the quarterly division of the period of the analysis.
Conclusions
Recent research on capital market issues are arranged into four dimensions: theory and assumptions of efficient capital market, government's role in it, other distortions and global interrelatedness. Main conclusions are decentralized by topics and summarized:
Effective capital market
− The effective market model's interpretations' evolution, from couple perfect market conditions to information symmetry as price forming constraint, was formed.
− The capital market effectiveness could be interpreted through several approaches: classical hypothesis of effective markets, information that forms price model, macroeconomic access as well as behavior of its participants (e.g. investors).
− Highly correlated issues of the capital market theory: reaching the perfection while eliminating the imperfection of information shared in the market.
− There exists the transition of effectiveness through related markets as well as their constraints.
Information (a)symmetries
− While watching positive correlations between macroeconomic indicators and information symmetry, the main causes of market imperfection are found in microeconomics.
− Effective capital market hypothesis emphasizes the correction of prices, and is explained as the symmetric and right information, possessed by market players.
− Symmetric information correctly allocates the resources, regarding the main capital market function of distribution.
Government's role
− There is an absence of impact of governmental restrictions on overall market development while influence on primary markets (e.g. capital) is being identified. Therefore the conclusion has come of action based rebound on action targeted market.
− The view of financial market dynamics can be expressed as follows: if borrowers and lenders are willing and able to contract, and liquidity providers find conditions conducive to trading the instruments that are created, then financial markets will develop. The regulatory structure can support this process by removing obstacles that make potential borrowers, lenders, and liquidity providers unwilling or unable to play their roles, and by creating the right incentives for each agent to fulfill their end of the bargain.
Global interrelatedness
− More disadvantages of foreign invasion by examining the causes, nature and impact of rising participation of foreign investor in local currency bond markets of developing country are seen.
− U.S. financial shocks reverberate around the world much more strongly than shocks from other regions, including the Euro area, while inward spillovers to the U.S. from elsewhere are minimal.
− The co-movement in emerging market bond returns and disentangles are of influence of external and domestic factors.
− Spillovers all across the world are identified and the impact to local capital markets is agreed.
The theoretical "eight" -dispersion of research topics is visualized with the division into dimensions of all recent research works by international organs (IMF, WB, OECD, different institutes and universities).
The definition of life-cycle of efficient capital market, which could be reached by diminishing the impact of such obstacles as information asymmetries, including government role with acceptable restrictions and getting interconnected in global market, is being introduced.
New approach to the classification of countries by their maturity in capital market is proposed: under-developed, emerging, integrated, matured. An argumentation comes from the agreement on opinion that country's economic conditions are determined by the behavior of financial markets, miserable statistics, lacking attention on global research and data of some countries, sequencing approach and logic for the existence of intermediate period from performing in emerging market rating to becoming a matured one. Forming the sequences of the countries by their capital market development stages more intuitively, the range scale is being of the important need to be implemented for overall, coherent and integral research and data collection.
Describing the state of art of bond market, in 2009-2012 there is no significant gain in development as well as sharp decrease of debt securities (bond market), slightly distinguishing the trend of matured capital markets. However, the mismatch in amounts of total debt securities and gross government debt is well-seen. | 7,856.8 | 2013-09-13T00:00:00.000 | [
"Economics"
] |
A Pre-Service Teacher Experiences of Creating Vocabulary Quizzes for EFL Adult Learners: the ACTIONS Model
IDLE (Informal Digital Learning of English) is a worldwide phenomenon that represents one of the most significant advances in autonomous language learning outside the classroom in recent decades. This study examines the experiences of IDLE activities based on the ACTIONS model (Access, Cost, Teaching and learning, Interactivity and user-friendliness, Organizational issues, Novelty, and Speed) which focused on vocabulary. The results of the study are intended to be a self-reflection on the factors involved in creating English vocabulary quizzes on Instagram as IDLE sources for higher education students. The study aims to use social media, especially Instagram, as a learning tool in a digital context. The researcher uses written narratives that contain her experiences in creating such English vocabulary quizzes. For that reason, the study participant is the researcher herself, a 21-year-old female undergraduate student in the English Education Department. In the study, the researcher uses thematic analysis to analyse the narrative data. This includes three activities: 1) repeated reading of the data; 2) coding and categorising the data extracts; and 3) recognising the thematic headings. The results show that creating IDLE activities based on the ACTIONS model leads to flexibility of access, affordable costs, teaching and learning needs based on followers' feedback and correction, excellent interactivity and user-friendliness, no organisational issues, novelty, and speed. The study offers new insights into how English pre-service teachers' engagement with IDLE serves as a significant factor in their future teaching tasks.
INTRODUCTION
In this era, students have various ways of learning. Technology has rapidly spread all over the world, so most people around the world have access to digital devices. The new digital pedagogical approach to learning knowledge and skills is informal. Lee (2017) refers to this as Informal Digital Learning of English (IDLE). Lee and Lee (2018) explain that IDLE is self-directed, naturalistic, digital learning of English in unstructured, out-of-class environments and is independent of formal language settings. In other research, Lee and Dressman (2017) state that the concept of IDLE is based on autonomous learning, using many different digital devices (e.g., smartphones, Play Stations, MP3, TV, PC) and resources (e.g., blogs, YouTube, Twitter, Instagram, online games, etc.). This definition of IDLE suits students' daily life. Digital devices and resources are tools for practising language outside the classroom.
English is the most widely used worldwide language as a result of globalisation. As a result, many people nowadays consider strong English skills to be essential. The rapid proliferation of wireless networks and mobile devices has made English learning in a mobile learning context popular in recent years. Many studies have developed mobile English learning systems to aid the study of English vocabulary any time and from any location (Chang & Hsu, 2011;Chen & Chung, 2008;Chen & Hsu, 2008;Chen & Li, 2010;Chen & Tsai, 2010;Cheng, Hwang, Wu, Shadiev, & Xie, 2010;Oberg & Daniels, 2013). Technology has aided the development of learning tasks, particularly language learning, easier. Social media is one of technology's support components. According to Dewing (2010), it is a broad category of web-based and mobile services that enable users to participate in online conversations, contribute to user-generated content, and form online communities.
Throughout the past decade, CALL research has presented data on the relationship between L2 vocabulary learning and extramural language learning over mobile devices (Stockwell, 2010 Hayati, Jalilifar, & Mashhadi, 2013;Lu, 2008) and social media (Sockett & Toffoli, 2012). IDLE research on digital games and L2 vocabulary acquisition has been proliferating in European countries. As found by Olsson (2011), Swedish teens who most frequently experienced out-of-school extramural English (EE) activities (mostly digital games) achieved the highest scores in their English writing assignments and used deeper and better English vocabulary. This finding is in line with other studies (Sundqvist, 2009;Sylvén & Sundqvist, 2012), which show that repeated EE activities (mostly digital games) are closely associated with better L2 English vocabulary acquisition. In addition, a recent study of IDLE in Indonesia was conducted by Lee and Drajati (2019). They found that IDLE practices and affective variables were significantly related to students' willingness to communicate. However, the term 'IDLE' is currently in its early stages. As a result, there have been few related studies, particularly ones focusing on skills or practice. This means that no study of extramural English activities related to quizzes on Instagram as IDLE sources for higher education students has been conducted to date. Based on this gap, such study needs to be conducted in this digital era in Indonesian society. Therefore, this work explores a pre-service teacher's experiences of creating extramural English vocabulary quizzes on Instagram as IDLE learning sources.
THEORETICAL FRAMEWORK ACTIONS Model for Developing English Vocabulary Quizzes on Instagram as IDLE Sources
Young people can now make good friends using technology, and take part in socialising in a digital world (Ito, M., Horst, H. A., Bittanti, M., Herr Stephenson, B., Lange, P. G., Pascoe, C. J., & Robinson, L.,2009). Over 90% of Americans aged 18-29 possess mobile devices such as smartphones, laptops, tablet computers, etc. and access social networking sites (Rainie, 2016). This has also significantly influenced how students learn (Prensky, 2001;Bennett, Maton, & Kervin, 2008). Junco (2012) states that social media plays a role in improving higher education students' academic results.
Designing, acquiring, and implementing technology for teaching and learning is one of the most challenging issues in higher education today. The rapid advance of technology has forced educators to stay abreast of developments. However, technology can be either good or bad, depending on how it is used. Bates (1995) proposes the ACTIONS model as a decision-making framework for evaluating technology in terms of learner accessibility, cost structure, teaching and learning application, interactivity and user-friendliness, organisational impact on the educational institution, novelty, and the speed with which courses can be developed for the technology. ACTIONS encompasses a set of critical questions that teachers need to ask before using certain technology for education and training purposes. These aspects are helpful in analysing whether or not teachers should experiment with a particular technology. For example, Bates (1995) suggests that technology-driven media can enable flexible, career-oriented learning in various locations, including the office, study groups, and the home. This is related to the IDLE (Informal Digital Learning of English) concept, which is flexible and can be undertaken in various locations. Previous research (Lee, 2018;Lee & Kim, 2014;Sockett & Toffoli, 2012) has suggested that IDLE activities can take place outside of the classroom, in media-rich environments such as the home, public transportation (e.g., bus and metro), and café or restaurant settings. In other words, by participating in IDLE activities, EFL students could maximise their learning chances outside of the language classroom. Although this aspect of CALL inquiry is still at its initial stage, the recent literature presents mixed results on the effectiveness of IDLE in language learning. Supported theoretically by concepts of incidental language learning (Schmidt, 1994), learner autonomy (Holec, 1981), and informal language learning (Benson, 2011), IDLE (Lee, 2017) can be theorised as self-directed, informal English learning suitable for a variety of different digital devices (e.g., smartphones, PCs, tablet computers) and resources (e.g., websites, social media). However, Burston (2014Burston ( , 2015 and Sung, Chang, and Yang (2015) claim there are no significant effects of mobile language learning in informal settings. On the other hand, other studies found a positive relationship between IDLE and L2 outcomes, such as in reading and listening (Sylvén & Sundqvist, 2012), speaking (Mitra, Tooley, Inamdar, & Dixon, 2003), writing, (Sun et al., 2017), vocabulary (Jensen, 2017;Sundqvist & Wikström, 2015;Sylvén & Sundqvist, 2012), and formal testing (Lai, Zhu, & Gong, 2015;Sundqvist & Wikström, 2015).
PAPER
IDLE is categorised into two domains: extracurricular and extramural c. The former implies a semiautonomous L2 activity in an out-of-class digital setting that a formal language teacher still organises. In other words, it refers to students' willing autonomous English learning in an out-of-class digital context, which is still assessed by a teacher (Lee, 2019). The latter refers to an autonomous L2 activity in an out-of-class digital setting which is unrelated to formal language instruction (Lee, 2019). It is different to IDLE in the extracurricular context in that the teacher is not involved in the students' behaviour. In other words, it implies that EFL students are taking part in autonomous English learning in digital, unstructured, out-of-class settings that are not allied to a formal institution and not assessed by a teacher. For example, L2 students may independently join Facebook, Instagram, or Twitter comment sections in English to beinvolved with other learners.
This study focuses on IDLE in the extramural context of Instagram development of English vocabulary quizzes. As stated previously, this idea comes from the researcher's experiences of creating such quizzes on Instagram. Initially, this was just for fun, until it began to attract positive feedback from her followers. Therefore, the researcher began to create English vocabulary quizzes for her followers as a learning resource. In addition, the researcher believes that vocabulary is a significant aspect of language acquisition. According to Nation (2013) and Willis and Ohashi (2012), vocabulary is a fundamental component of any language, and so is a vital element of second language (L2) acquisition.
Vocabulary knowledge affects both receptive skills (reading and listening) and productive ones (speaking and writing) and shows a critical success of general language proficiency (Alderson, 2007;Laufer & Goldstein, 2004). L2 learners admit that a lack of vocabulary causes them difficulties in acquiring, comprehending, and using the language (Nation, 2013). According to González-Fernández and Schmitt (2020), vocabulary acquisition refers to all the processes concerned with learning lexical items (i.e., single words and formulaic language) in great depth in order to use them both productively and receptively help of multiple incidental and intentional encounters with these items in various contexts.
Studies have observed that L2 students' vocabulary knowledge embraces familiarity with a series of words, obtaining many types of knowledge about each word, and finding connections between multiple lexical items to build semantic networks (Cremer, Dingshoff, Beer, & Schoonen, 2010). However, it is still not clear how vocabulary is stored and processed in the lexicon. Furthermore, it is acknowledged that all words are interconnected in multiple ways, so learning one lexical item influences the learning of others (Meara & Wolter, 2004). Therefore, to have complete knowledge of a word, making a rich and interrelated mental lexicon that supports more rapid, comprehensive, precise networks among words (Cremer et al., 2010). Nevertheless, investigating the connections between words is a complex and challenging task.
Since the knowledge of words is acquired through many language experiences, such as explicit instruction and incidental exposure (Schmitt, 2008), their acquisition is not a static process. More accurately, word knowledge is a dynamic process that constantly changes and develops, meaning that the acquisition of a word goes through different phases until all the related elements (such as form-meaning mapping, collocational information, and word parts) are understood (Fitzpatrick, 2012).
Based on the definitions above, vocabulary can be defined as one of the essential components of language proficiency, including receptive skills (listening and reading) and productive ones (speaking and writing). Vocabulary acquisition is about learning single words, understanding each word, and connecting multiple lexical items to semantically build a language system. The more incidental and intentional encounters with lexical items in various contexts, the more vocabulary knowledge learners will have. 3. METHOD This is qualitative research using a narrative inquiry. The study participant is the researcher herself, a 21year-old female undergraduate student in the English Education Department. In the study, she used written narratives containing her experiences in creating English vocabulary quizzes on Instagram as IDLE sources. As stated by Barkhuizen, Benson and Chik (2014), some scholars have highlighted that the main strength of narrative inquiry lies in its focus on how people use stories to understand their experiences in the areas of inquiry in which it is essential to understand phenomena from the perspectives of those who experience them. The purpose is to attain a detailed story of experiences and understanding of the entity (the theme). In narrative research, there are two common approaches: biographical and autobiographical.
PAPER
In this study, the researcher used an autobiographical approach, which meant she analysed or told her own stories in three diaries. These comprised the primary data for this research. The diaries were written based on the quiz topic. They consisted of the past, present, and future of creating English quizzes on Instagram. To support the primary data, the researcher provided screenshots as artefacts.
Thematic analysis was used to analyse the narrative data, as suggested by Barkhuizen, Benson and Chik (2014). It includes three activities: 1) repeatedly reading the data; 2) coding and categorising the data extracts; and 3) recognising the thematic headings. Good thematic analysis always requires repeated reading of the data and several rounds of analysis, in which the researcher goes back and forth between the data and its coded and categorised forms to perfect themes and theoretical relationships. The data were analysed based on the ACTIONS model and IDLE research instrument.
RESULTS AND DISCUSSION RESULTS
This chapter presents the research results and discussion. The themes in this section result from careful coding of the data set by thematic analysis, which was intended to explore the stories of the researcher's extramural English activities in a digital context (social media). The themes below are based on the framework proposed by Bates (1995) in relation to the ACTIONS model.
Access
This concerns how accessible a particular technology is for learners and its flexibility for a particular target group. The findings show that creating an Instagram account is easy and free. No subscription is needed and it allows flexibility of access, either in or outside of class, using a smartphone (virtually anytime, anywhere
Cost
Teachers should consider the unit cost per student when they wish to try a particular technology. Instagram is still affordable to obtain and run; it is just necessary to download it from the Play Store for Android or the App Store for iOS. It is good that downloading Instagram is free, although a data plan or Wi-Fi connection is needed to run the app. Creating
Teaching and learning
This section discusses what kinds of learning are needed, what instructional approaches will best meet these needs and the best technologies for supporting such teaching and learning. For example, Instagram can support teaching and learning needs from followers' requests and feedback. Another finding is that followers' correction allowed the researcher to learn and recheck before posting the quizzes.
There is some positive feedback from my followers, such as 'Please make more quizzes like these,' 'Go on,' 'more more more,' 'Thank you, I learned a lot from these quizzes even though I still made mistakes, 'Please make quizzes before exams,' 'Create more content like these, pretty good for learning English,' and many more.
(Pre-service teacher, 2019, in a boarding house)
I typed the word 'weather' in the wrong order 'waether'. Then, one of my followers, HL (accronym for a follower's username in Instagram -Editor), slid into my Direct Message (DM) and told me there was a typo in my quiz, which might happen because I was in a rush while typing the quiz. That was true. I typed the items quickly because I typed them one by one. However, I thought I had to recheck what I had typed to make sure there were no typos next time. Another follower, NT, slid into my DM, telling me I mistyped a word. The item showed a past tense sentence, but I wrote the multiple choices using the present tense. For example, the answer was 'saw,' but I wrote all the multiple choices in the present tense 'see,' 'watch,' 'look.' I realized that she was right, and that I should be more careful next time. She also said that it might be because I was in a hurry.
(Pre-service teacher, HL & NT, 2020, in a boarding house, on Instagram) The word weather is a typo, isn't it? Fig. 1.8: It is not seen but saw. I know that you know, but perhaps you were in a hurry
Interactivity and user-friendliness
This section explains what kind of interaction Instagram enables and how easy it is to use. It contains features to make English quizzes, such as a quiz sticker, question sticker, poll sticker, and pictures. As the quiz creator, the researcher can see each follower's answers to the quiz. Moreover, Instagram has direct messaging to receive followers' questions and respond to their answers.
Instagram has many features that can be used for creating English quizzes. The researcher has tried some of these, such as a quiz sticker, question sticker, and pictures. For the quiz sticker, followers just need to click on the choices provided. Then, they will directly know whether they get the correct answer or not. They will see a checkmark and green display when they have given the correct answer. When they get the wrong answer, they will see a cross mark, a red display, and key answer. For the question sticker, followers need to type the answer in the column provided. However, they cannot know whether they have given the right or wrong answer. There is no particular display on the quiz sticker. Therefore, I respond to their answers via DM one by one.
Organisational issues
This section deals with the organisational barriers that may need to be faced before a particular technology can be used successfully. In fact, no barriers were found before using Instagram to develop English quizzes. This can be seen in the following diary excerpt.
I didn't have any barriers before using Instagram to develop English quizzes. Also, I am not connected to a particular organization or institution in making English vocabulary quizzes on Instagram. In other words, no individual forces me to do it. Moreover, the ease of access and interactivity also influences the flexibility of using Instagram.
(Pre-service teacher, 2019, in a boarding house)
Novelty
Instagram was launched in October 2010. Therefore, it is not considered a brand new technology. However, it also means that it should be easy to use for people worldwide.
Instagram was launched in October 2010. It has been around for some time and has a large user base. So, it ought to be available and replicable around the world. In other words, everyone can easily have Instagram for any purpose. Therefore, people will not find it challenging to use it.
(Pre-service teacher, 2019, in a boarding house)
Speed
This is about how quickly this technology can mount and how quickly materials can be changed. Instagram is able to include stories by adding a brief introduction, posting the quizzes, and saving the highlights. Instagram highlights are beneficial because stories expire after 24 hours. Therefore, highlights are helpful when the researcher and followers want to look back at the quizzes at any time.
It is speedy. I usually added a brief introduction about a specific topic before jumping right into the quizzes. Once I was ready with the intro, I straightaway posted it on my Instagram story. The next step was to post the quizzes one by one. After all the quiz items were posted, they would expire after 24 hours.
DISCUSSION
This discussion section addresses the all theoretical and practical implications summarised from the findings. Social media has many impacts on students' learning experiences. It comprises platforms that can help student teachers learn and practise their active English; for example, providing English vocabulary quizzes on Instagram. The findings of this research on Bates' (1995) ACTIONS model indicate that Instagram can be a medium that is able support a learning process. Pre-service teachers and learners can easily access it if they already have an account. Those who do not have an Instagram account can create one using email or phone numbers or link it to their Facebook account. There is no subscription, so it is free to use. Instagram also allows flexibility of access, either in or out of class, using a smartphone (almost anytime, anywhere).
Because of the rapid advance sand expansion of digital technology capabilities, i tis now possible to use and generate content in English at any time and from any location, resulting in various learning opportunities (Lee & Lee, 2020). These possibilities can potentially change the "face of language learning" for the better (Richards, 2015). Today's EFL students are increasingly absorbing and utilising English informally through a variety of digital resources, including social media, virtual communities, language study apps, and large multiplayer online roleplaying games (MMORPGs) (Dressman & Sadler, 2019;Lai, 2017;Sauro & Zourou, 2019;Sockett, 2014;Sundqvist & Sylvén, 2016). This is related to the tenets of IDLE (Lee & Lee, 2020), which uses self-directed English activities in informal digital settings, motivated by personal interests and undertaken independently without being assessed by a teacher.
Another consideration is cost. It is relatively affordable to obtain and run the app. It is simply necessary to download it from Play Store for Android or the App Store for iOS. It is also free to download the App. To run the app, you need a data plan or wi-fi connection. Nowadays, most people own a smartphone and monthly data plan, so the cost is reasonable.
Regarding the third aspect of the ACTIONS model, which is teaching and learning, Instagram can support their needs from followers' requests or feedback. Their feedback shows that they tend to enjoy learning English through the platform. Moreover, followers' correction also allows the researcher to learn herself and recheck before posting the quizzes. Lee and Lee (2020) state that IDLE activities can be possible ways for students to have pleasant emotional experiences in English learning. Teachers can observe and note types of IDLE activities their students are already engaged in (e.g., watching YouTube clips in English or chatting with others via social media) and design and integrate such IDLE-embedded activities into in-class or out-of-class assignments because students primarily structure and implement IDLE activities.
Related to interactivity and user-friendliness, Instagram has features to make English quizzes, such as a quiz sticker, question sticker, poll sticker, and pictures. From the quizzes, the researcher, as the quiz creator, can see each follower's answers. To respond to these, a Direct Message can be used. Another Direct Message function allows the researcher to receive followers' questions about the quizzes. Several studies have examined the concept of learning or perceptions of learning arising from social media, including ones on technological competence (Dymoke & Hughes, 2009;Salminen et al., 2016) and academic skills (Dymoke & Hughes, 2009;Kiliç & Gӧkdaş, 2014;Kivunja, 2015;Wheeler & Wheeler, 2009). Social constructivism highlights the concept of learning or perceptions of learning which rely mainly on a sense of community (Kiliç & Gӧkdaş, 2014), in which participants are encouraged to collaborate, share, discuss, and challenge ideas and beliefs. The researcher has found Instagram to be useful for such activities by providing question stickers for sharing. This feature allows the researcher to learn English from her followers in a process of knowledge sharing. First, the researcher shares English quizzes. Then, she adds a question sticker for her followers to share their knowledge. Subsequently, the researcher responds to what has been shared by her followers. In responding, she continuously learns something new from her followers. Overall, the interactivity and user-friendliness of Instagram are excellent. Surprisingly, the researcher has not found any barriers to using Instagram for developing English quizzes. The ease of access and interactivity support the process. In addition, the researcher is not connected to a particular organisation or institution in developing English vocabulary quizzes on Instagram. In other words, she is not forced to do it. This means no organizational issues are involved. The next aspect concerns novelty. Instagram was launched in October 2010. Therefore, it is not considered a brand new technology, which means it should be easy to run for people worldwide. With regard to the final aspect of the ACTIONS model, speed, Instagram can include stories at one time by adding a brief introduction, posting the quiz items, and saving them on the highlights. However, Instagram stories will expire after 24 hours. That is why the researcher saves the quizzes on the highlights so that she and her followers can look back at the quizzes again anytime. However, Instagram is not yet perfect in comparison to other teaching and learning platforms, such as Kahoot, Quizizz, Moodle, Microsoft Teams, and others. Even so, Instagram can be an option to be used in or out of the classroom. Other factors include the impact on practice, as participants showed pre-service teachers would integrate similar tasks into their teaching in the future (Bravo & Young, 2011;Fisher & Kim, 2013).
CONCLUSION
This research shows that Instagram can be an option for teaching and learning in or out of the classroom. Based on the ACTIONS model, it is a platform that should be tried, although it is not a perfect educational platform, such as Moodle, Kahoot, Microsoft Teams, etc. However, there are some considerations as to why preservice teachers should experiment with Instagram to support their teaching and learning. First, it is a question of access. The creation of an account is easy and free, and the platform allows flexibility of access anytime, anywhere using a smartphone. This is also related to the concept of IDLE, which also underlines the flexibility of access. Second, the costs to run Instagram are relatively affordable. Third, in relation to teaching and learning, Instagram can help pre-service teachers fulfil followers' learning needs online. Moreover, it has several features that facilitate good interactivity. Therefore, the researcher has not found organisational issues when using Instagram for developing English quizzes. In terms of novelty, Instagram is not a brand new technology, so learners are already familiar with it. The final consideration is that Instagram allows quick inclusion of materials. This research implies that creating IDLE activities on Instagram for higher education students is an experience that can inform the researcher's future career. She believes that social media activities can be brought to the classroom, especially using Instagram. Social media activities can be used as a variation in learning in the classroom so that students do not become bored. Moreover, the researcher also considers social media, especially Instagram, to be a learning memoir for what has been learned. As a pre-service teacher, the researcher can practise her ability to produce English items, utilise Instagram as a learning platform, and use certain Instagram features that attract more followers.
This study suggests that teachers can encourage and promote the use of social media in their future careers. Therefore, higher education pre-service teacher courses should actively teach social media usage. The researcher also hopes that students can learn English autonomously by using social media, such as Instagram, Facebook, Twitter, and TikTok, among many others. Students can also practise and share their English knowledge using their social media. Finally, the study could help further research. Other researchers could use the results of this research, which present a narrative study of a pre-service teacher's perspectives, in future IDLE studies. In addition, future researchers could consider pre-service teachers' perspectives and the English context when conducting IDLE research. | 6,316.8 | 2022-04-30T00:00:00.000 | [
"Education",
"Computer Science"
] |
The lipopeptides pseudofactin II and surfactin effectively decrease Candida albicans adhesion and hydrophobicity
A serious problem for humans is the propensity of Candida albicans to adhere to various surfaces and its ability to form biofilms. Surfactants or biosurfactants can affect the cell surfaces of microorganisms and block their adhesion to different substrates. This study investigated adhesion of C. albicans strains differing in cell surface hydrophobicity (CSH) to polystyrene microplates in order to compare the ability of lipopeptide biosurfactants pseudofactin (PF II) and surfactin (SU) to prevent fungal adhesion to polystyrene. The biosurfactants decreased adhesion of tested strains by 35–90 % when microplates were conditioned before the addition of cells. A 80–90 % reduction of adhesion was observed when cells were incubated together with lipopeptides in microplates. When microplates were pre-coated with biosurfactants, PF II was less active than SU, but when cells were incubated together with biosurfactants, the activity of both compounds was similar, independent of the CSH of strains. When cells were preincubated with lipopeptides and then the compounds were washed out, the adhesion of hydrophobic strains increased two times in comparison to control samples. This suggests irreversible changes in the cell wall after the treatment with biosurfactants. CSH of hydrophobic strains decreased only by 20–60 % after incubation with biosurfactants while adhesion decreased by 80–90 %; the changes in cell adhesion can be thus only partially explained through the modification of CSH. Preincubation of C. albicans with biosurfactants caused extraction of cell wall proteins with molecular mass in the range of 10–40 kDa, which is one possible mechanism of action of the tested lipopeptides. Electronic supplementary material The online version of this article (doi:10.1007/s10482-015-0486-3) contains supplementary material, which is available to authorized users.
Introduction
Candida albicans is responsible for fungaemia, especially in immunocompromised patients. Cell features that cause mycoses encompass, e.g., adhesion, secretion of hydrolytic enzymes, filamentation and hydrophobicity (Verstrepen and Klis 2006). Understanding how C. albicans morphogenesis modulates the molecular composition of the fungal cell surface and interactions with biotic and abiotic surfaces is important, but still unclear.
The microbial adhesion results from specific interactions between cell surface structures and the surface of the substrate, or from non-specific interaction forces, including Brownian movement, van der Waals attraction, gravitational forces and surface electrostatic charges. One of the important factors is the hydrophobicity of cell surface (Krasowska and Sigler 2014).
Cell surface hydrophobicity (CSH) is connected with adhesion and pathogenic processes of C. albicans. Hydrophobic cells are more adherent than hydrophilic ones to epithelial and endothelial tissues as well as to abiotic surfaces (Glee et al. 2001;Hazen 2004).
Hydrophobicity of C. albicans cells alters in response to changes in environmental conditions (e.g. temperature, composition of medium) and growth phases (Hazen et al. 2001) and can be switched between hydrophilic and hydrophobic phenotypes (Masuoka and Hazen 1997). Hydrophilic cells have an elongated acid-labile mannan fraction in the cell wall and the length of this structure affects the folding of cell wall fibrils (Masuoka and Hazen 1999). Chaffin (2008) supposed that Csh1 protein influences the acid-labile mannan composition, because of differences between hydrophobic and hydrophilic cells in mannan fractions. Mannoproteins can therefore be potential targets for new antifungal drugs (Gow et al. 1999).
Biosurfactants such as lipopeptides are particularly interesting as antifungals because of their high surface activity and antibiotic potential. Several natural lipopeptides, e.g. echinocandins, block specific enzymatic reactions in the synthesis of cell wall components (e.g. b-1,3-glucan or chitin). Lipopeptides such as surfacin (SU), iturin and bafilomycin disturb the plasma membrane (Makovitzki et al. 2006). The adsorption of biosurfactant molecules on a surface was found to change its hydrophobicity, which might cause changes in the adhesion processes (Zhong et al. 2008;Singh et al. 2013).
Previously, we described the antiadhesive activity of the lipopeptide pseudofactin II (PF II), produced by Pseudomonas fluorescens BD5 (Janek et al. 2010) against several uropathogenic bacteria and C. albicans, and did not detect a significant impact on C. albicans growth (Janek et al. 2012).
PF II and SU are both cyclic lipopeptides. In the PF II molecule a palmitic acid is connected to hydrophilic ''head'' of eight uncharged amino acids (Janek et al. 2010), whereas SU is a lipoheptapeptide linked to a b-hydroxyl fatty acid. Commercially available SU (Sigma-Aldrich) is a mixture of congeners that differ in the length of the carbon chain (C 12 -C 16 ). Moreover SU is negatively-charged because of Asp and Glu amino acids within the molecule (Raaijmakers et al. 2010). These differences cause a variations in the biological activity of lipopeptides e.g. disruption of plasma membrane by SU.
In this work we compared the action of PF II and SU on C. albicans strains that differ in CSH. We examined the influence of lipopeptides on the viability and adhesion of C. albicans on polystyrene. We also checked the impact of the biosurfactants on CSH of C. albicans. Our results suggest differences in the mechanisms of action between PF II and SU. Micelles of PF II and SU cause irreversible changes in the cell wall of hydrophobic strains of C. albicans but a decrease in adhesion could be explained only partially by the influence of lipopeptides on CSH. Moreover, the biosurfactants appeared to be able to extract some cell surface-associated proteins from C. albicans cell wall (CWP), which is demonstrated for the first time in this work.
Microorganisms and culture conditions
Biosurfactant-producing strain P. fluorescens BD5, obtained from freshwater from Arctic Archipelago of Svalbard, was cultivated in LB medium as described earlier (Janek et al. 2010). C. albicans strains (Table 1) were a generous gift from D. Sanglard (Lausanne, Switzerland) and were cultivated in 5 ml YPG broth containing 10 g/l bactopeptone (Difco, USA), 10 g/l yeast extract (Difco, USA), and 20 g/l glucose (Bioshop, Canada). Candida cultures were incubated at 28°C for 24 h without agitation and then stored at 4°C for a maximum of 2 weeks. All experiments were carried out on fresh C. albicans pre-cultures (4.85 ml of YPG inoculated with 150 ll of C. albicans culture and incubated for 24 h at 28°C). Before conducting the experiments C. albicans cells were centrifuged twice (10009g) for washing out the culturing medium and resuspended in PBS pH = 7.4 (8 g/l NaCl, 1.4 g/l Na 2 HPO 4 , 0.25 g/l KH 2 PO 4 , 0.2 g/l KCl) or phosphate buffer (PB; 16.9 g/l K 2 HPO 4 , 7.3 g/l KH 2 PO 4 ).
Production, isolation and purification of pseudofactin II (PFII)
For the production of PFII, P. fluorescens BD5 was cultivated in mineral salt medium (MSM) containing 7 g/l K 2 HPO 4 , 2 g/l KH 2 PO 4 , 1 g/l (NH 4 ) 2 SO 4 , 0.5 g/l sodium citrate 9 2H 2 O, and 0.1 g/l MgSO 4 9 7H 2 O supplemented with 20 g/l glucose at 28°C without agitation as described earlier (Janek et al. 2010). Briefly, 0.5 l of MSM was inoculated with 5 ml of P. fluorescens BD5 culture in LB (24 h, 28°C) and incubated for 1 week at 28°C without agitation. Cellfree supernatant was afterwards extracted three times with ethyl acetate. The solvent was evaporated under vacuum and crude extract was dissolved in methanol and purified by RP-HPLC (Janek et al. 2010).
Biosurfactant concentrations
Biosurfactants were tested in the final concentrations of 0.035 or 0.1 mg/ml for PFII and 0.005 or 0.015 mg/ ml for SU. These concentrations were chosen to test the influence of biosurfactant monomers (*0.5 9 CMC) and micelles (*1.5 9 CMC). PFII was extracted and purified as described above. SU was manufactured by Sigma-Aldrich (USA). Biosurfactant stock solutions were dissolved in PBS and stored at -20°C.
Antifungal activity of biosurfactants
The antifungal activity of biosurfactants was tested in 96-well flat-bottom polystyrene microplates (Sarstedt, Germany). We added 50 ll of double strength YPG and 50 ll of biosurfactant solution in PBS to each well or PBS to control wells. Every well was afterwards inoculated with overnight Candida culture in YPG to reach the initial optical density at 600 nm (OD 600 ) of 0.01. The microplates were then incubated for 24 h at 30°C. After incubation the OD 600 was measured with UMV 340 microplate reader (Asys Hitech, Austria). Antifungal activity of biosurfactants is expressed as a where OD T is the OD 600 of wells containing biosurfactants in PBS and OD C is the OD 600 of control samples (wells without biosurfactants).
Cell surface hydrophobicity (CSH)
For determining the effect of biosurfactants on C. albicans CSH, cell suspensions in PB were transferred to Eppendorf test tubes and PFII or SU stock solutions in PBS were added to reach the biosurfactant final concentrations. The same amount of PBS was added to the control samples. Suspensions were incubated for 2 h at 37°C with agitation (300 rpm) and then diluted to an OD 600 of 0.5. The MATH (microbial adhesion to hydrocarbon) was used to evaluate the CSH of Candida cells (Coimbra et al. 2009). Briefly, 2 ml of the cell suspension in PB were moved to a glass tube (100 9 15.5 mm) and 100 ll of hexadecane w added. The samples were then vortex-shaken for 3 min and the phases were allowed to separate for 1 h. The OD 600 of the aqueous phase was measured and CSH, defined as percentage of cells adhering to hexadecane, was calculated as follows: where OD 600 is the optical density of the aqueous phase at 600 nm. In modified trials, biosurfactants were washed out (centrifugation 10009g) with PB before diluting cell suspensions to an OD of 0.5 and measuring CSH. Adhesion of Candida albicans to polystyrene PF II and SU were tested as C. albicans adhesioninhibiting agents in flat-bottom 96-well polystyrene microplates (Sarstedt, Germany) in three different assays. In pre-adhesion assay, microplate wells were preincubated with 100 ll of biosurfactant solutions in PBS for 2 h at 37°C with agitation (300 rpm). PBS buffer was used as a positive control. Subsequently, wells were washed two times with PBS. C. albicans suspensions in PBS were diluted to give an OD 600 of 0.6. The highest adhesion of C. albicans strains to polystyrene was observed at this OD (Janek et al. 2012). 100 ll of Candida suspensions were added to wells and incubated for 2 h at 37°C with agitation (300 rpm). Then supernatants were removed and wells were washed two times with PBS to remove nonadherent cells. Adherent cells were stained with 0.1 % crystal violet for 5 min and then wells were washed three times with PBS. The dye was released by 200 ll of 0.05 M HCl with 1 % SDS in isopropanol and the absorbance at 590 nm (Abs 590 ) was read off with Asys UMV 340 microplate reader (Asys Hitech, Austria). Cell adhesion was expressed as the Abs 590 or as the percentage of Abs 590 of control samples (100 %): where Abs t is the Abs 590 of wells pretreated with biosurfactants and Abs c is the Abs 590 of control wells (pretreated with PBS only). In addition, we tested C. albicans adhesion to microplates in the presence of biosurfactants. Briefly, we added biosurfactants to Candida suspensions in PBS to reach final concentrations and the OD 600 of 0.6. The same amount of PBS was added to the control samples. Then, 100 ll of suspensions were added to microplate wells and incubated for 2 h at 37°C with agitation (300 rpm). The microplates were washed, stained and read as described before. We also investigated the influence of preincubation of C. albicans strains with biosurfactants on their adhesion abilities. In brief, Candida cell suspensions in PBS were transferred to Eppendorf test tubes and biosurfactants were added to the desired final concentrations. The same amount of PBS was added to the control samples. Suspensions were incubated for 2 h at 37°C with agitation (300 rpm) and diluted to an OD 600 of 0.6. Then, 100 ll of suspensions were added to microplate wells and incubated for 2 h at 37°C with agitation (300 rpm). Microplates were washed, stained and read as described before. In modified trials, biosurfactants were washed out (centrifugation 10009g) with PBS before diluting cell suspensions to an OD 600 of 0.6 and conducting the adhesion assay.
Extraction of cell-wall associated proteins (CWP) by biosurfactants
We also tested if the addition of biosurfactants can cause extraction of proteins from the C. albicans cell surface. To conduct the experiment, Candida cell suspensions in PBS were transferred to Eppendorf test tubes and biosurfactants were added to the final concentrations. The same amount of PBS was added to the control samples. Suspensions were incubated for 2 h at 37°C with agitation (300 rpm). Then cells were removed by centrifugation (10009g) and filtration (0.2 lm). Proteins in supernatants were concentrated with Amicon Ultra 0.5 mL 3 kDa centrifugal filters (Millipore, USA). Concentrated samples were mixed with 96 denaturation buffer (150 mM Tris; 0.6 M EDTA; 12 % SDS; 60 mM DTT), heated at 95°C for 5 min and loaded onto 15 % polyacrylamide gel. Silver-stained gels were photographed with Chemi-Doc System (Bio-Rad, USA).
Fluorescence microscopy
Candida cell suspensions in PBS were transferred to Eppendorf test tubes and biosurfactants were added to the final concentrations. The same amount of PBS was added to the control samples. SDS was added to the final concentration of 1 % and served as positivecontrol samples. Suspensions were incubated for 2 h at 37°C with agitation (300 rpm) as described above. Then cells were centrifuged twice (10009g) and resuspended in PBS buffer. PI from stock solution (Bioshop, Canada) was added to the final concentration of 6 lM and suspensions were incubated for 5 min at room temperature. Next, Candida cells were pelleted and washed twice with PBS. 4 ll of Candida pellets were viewed with Zeiss Axio Imager A2 fluorescence microscope. All described assays were carried out at least three times in three replicates. Statistical analyses were performed using paired t test with Bonferroni correction. P values of \0.05 were considered significant.
Results and discussion C. albicans can use various carbon sources: glucose, galactose, fructose or hydrocarbons. Carbon sources at different concentrations promote changes in the structure of cell wall (McCourtie and Douglas 1981); thus increasing sugar concentration in the medium from 50 to 500 mM resulted in the production of an outer fibrillar-floccular layer of mannoproteins and also a linear increase of adherence to acrylic surfaces (McCourtie and Douglas 1981). Different culture conditions have therefore an impact on surface properties of Candida cells (Hobden et al. 1995). In our collection of C. albicans strains (Table 1), CAF4-2 and DSY653 were more hydrophobic than other strains (P \ 0.001) (Fig. 1).
A change in glucose concentration in the medium from 2 to 0.2 % decreased CSH but only in the case of two strains with the highest hydrophobicity (Fig. 1). These results suggest differences in cell wall composition and metabolism of URA3 mutants as reported earlier (Bain et al. 2001). Our results also indicate an impact of the site of integration of URA3 in C. albicans genome on changes in surface properties. Strains DSY653 and DSY1050 that vary in the site of integration of URA3 differ in some aspects such as CSH (Fig. 1).
Microbial surfactants often have antimicrobial properties but knowledge about mechanisms of their action is scarce. A few studies have shown that rhamnolipids increase the membrane permeability and alter its barrier function, causing cell damage (Sotirova et al. 2008). Lipopeptides such as SU, iturin or lichenisyn form ion-conducting membrane channels (Pueyo et al. 2009;Bensaci et al. 2011). In contrast to many other lipopeptides (Peypoux et al. 1999;Grangemard et al. 2001), PF II showed much weaker antimicrobial activity against bacterial and C. albicans strains (Janek et al. 2012). Also SU in tested concentrations exhibited no antifungal activity (Fig. 2). PF II was found to possess an antiadhesive, concentration-dependent activity against bacteria and yeast. The highest reduction of adhesion (80-99 %) was observed for C. albicans wild-type strain SC5314 (Janek et al. 2012). PF II was effective above the critical micelle concentration (0.072 mg/ml) and the adhesion was thus inhibited more strongly by micelles than by monomers (Janek et al. 2012). The microbial adhesion depends on the composition of the outer cell layer and is connected with hydrophobic/hydrophilic and ionic properties of the cell as well as with the properties of the polystyrene surface of microplates used in experiments (Neu 1996). PF II, due to its nonionic character, can probably coat positively or negatively charged surfaces, changing their properties.
We studied the adhesion of C. albicans to polystyrene microplates in a number of different experiments to compare the ability of PF II and SU to prevent fungal adhesion to abiotic surfaces. It is obvious that strains CAF4-2 and DSY653 have modified surface properties, but the nature of these changes is not clear (Bain et al. 2001).
We observed a decrease in adhesion of all tested C. albicans strains when the microplates were pretreated with PF II before the addition of the microorganisms (pre-adhesion assay) (Fig. 3). PF II was more active in concentrations higher than CMC (0.1 mg/ml) (Fig. 3). We observed a similar concentration-dependent effect for SU used as a standard lipopeptide biosurfactant, which decreased the adhesion even more than PF II (P \ 0.001) (Fig. 3). CAF4-2 and DSY653 adhered to the polystyrene microplate surface better than the other strains (P \ 0.01) and were able to adhere to a surface pretreated with lipopeptides more strongly than other strains (P \ 0.001) (Fig. 3).
Surprisingly, when cells and lipopeptides were incubated together for 2 h in the polystyrene microplate, the adhesion was blocked even more strongly (Fig. 4). Both PF II and SU micelles reduced C. albicans adhesion by *90 %. As for biosurfactant monomers, the action of lipopeptides was different. In this case, PF II was found to be a better antiadhesive agent than SU (Fig. 4). The antiadhesive activity of SU was similar to the situation when it coated the microplate before the addition of Candida suspension (cf. Figs. 3, 4). PF II was less active than SU in the case of hydrophobic strains when the microplate was coated before the addition of cells but when hydrophobic cells were incubated together with PF II, their adhesion decreased like in hydrophilic strains (Figs. 3,4). These results suggest differences in the mechanisms of action between PF II and SU, e.g. interactions between cell surface and/or polystyrene.
Interesting results were observed when the cells were preincubated with biosurfactants and the adhesion of coated and non-coated cells to polystyrene microplate was investigated (Fig. 5). When present in the solution (Fig. 5a), lipopeptides act as strong antiadhesives in micellar concentrations. PF II monomers reduced the adhesion of hydrophilic strains approximately two times and did not alter adhesion of hydrophobic strains CAF4-2 and DSY653 (Fig. 5a). Monomers of SU did not change adhesion of hydrophilic strains and increased it in the case of hydrophobic strains (Fig. 5a). During incubation of Candida cells with the biosurfactants, the predisposition of cells to adhesion changed and was different from the case when the microplate was pre-coated with PF II or SU (Figs. 3,5). However, micelles of PF II decreased adhesion to the same low level (10-20 %) as in experiments with a 2-h adhesion of cells coated with PF II (Fig. 4). When the biosurfactants were washed out before conducting the experiment, the adhesion of hydrophilic strains was comparable to control samples whereas for hydrophobic strains adhesion increased approximately two times (Fig. 5b). This result suggests irreversible changes in the cell wall of hydrophobic strains of C. albicans caused by micelles of PF II and SU after a 2-h incubation.
The microbial ability of adhering to different surfaces is connected with CSH, hence our intention was to investigate the influence of lipopeptides on Candida CSH. Biosurfactants can change CSH due to adsorbing to the cell surface (Kaczorek et al. 2013), like rhamnolipids which strongly adsorbed on the cell surface of yeast (Kaczorek et al. 2008).
After a 2-h incubation with PF II or SU, CSH of C. albicans CAF4-2 and DSY653 significantly decreased and this effect was concentration-dependent. Monomers of PF II influenced CAF4-2 and DSY653 more strongly than monomers of SU. Other tested strains seemed resistant to the influence of biosurfactants (Fig. 6a). On the other hand, when biosurfactants were washed out, CSH level of hydrophobic cells recovered (Fig. 6b). In this assay the time of incubation with biosurfactants was 2 h and these conditions can be compared to experiments with adhesion of cells treated with biosurfactants (Fig. 4). CSH of hydrophobic strains decreased only by 20-60 % (Fig. 6) while adhesion decreased by 80-90 % (Fig. 4). Also the potential irreversible changes in the cell surface of C. albicans caused by lipopeptides have an impact on adhesion but not on CSH of hydrophobic strains (cf. Figs. 5, 6). This result suggests that decrease in cell adhesion by lipopeptides can be only partially explained by the modification of CSH and should be considered only in the case of hydrophobic strains CAF4-2 and DSY653.
One of the mechanisms of action of lipopeptides on C. albicans cells could be a decrease in the level of some compounds (e.g. chitin, b-1,3-glucan) in the cell wall (Bizerra et al. 2011). Some protocols for the fractionation of fungal cell walls include treatment with synthetic surfactants (Pitarch et al. 2002;Klis et al. 2007). Therefore, we isolated several proteins from cellfree supernatants after preincubation of C. albicans cells with biosurfactants and visualized them on silverstained polyacrylamide gels (Fig. 7). We determined molecular masses of these proteins after SDS-PAGE electrophoresis to be in the range from *10 to 40 kDa and observed no differences between the action of PF II and SU or between hydrophobic and hydrophilic strains (Fig. 7). Simultaneously, PAS (Periodic acid-Schiff) staining for glycoproteins showed no bands on the gels (data not shown). Therefore, partial disruption of cell wall and extraction of cell surface-associated proteins can be the possible mechanism of the action of lipopeptide biosurfactants on C. albicans.
To exclude the possibility of contamination of cellfree supernatants (Fig. 7) with cytoplasmic proteins, we analyzed viability and membrane permeability of Candida cells with fluorescence microscopy (Fig. 8). The lack of propidium iodine (PI) fluorescence in control samples and cells incubated with lipopeptides indicate that cells were viable and membranes permeability was undisturbed (Fig. 8), which also confirms viability results shown earlier (Fig. 2). In contrast, cells treated with 1 % SDS showed significant fluorescence of dead cells.
Open Access This article is distributed under the terms of the Creative Commons Attribution 4.0 International License (http:// creativecommons.org/licenses/by/4.0/), which permits unrestricted use, distribution, and reproduction in any medium, provided you give appropriate credit to the original author(s) and the source, provide a link to the Creative Commons license, and indicate if changes were made. | 5,329.8 | 2015-05-29T00:00:00.000 | [
"Biology",
"Chemistry",
"Environmental Science",
"Medicine"
] |
Electron electric dipole moment and electroweak baryogenesis in a complex singlet extension of the Standard Model with degenerate scalars
We study the possibility of electroweak baryogenesis in the standard model with a complex scalar field, focusing mainly on a degenerate scalar scenario. In our setup, CP violation is provided by dimensional-5 Yukawa interactions involving the complex scalar field. In contrast to previous studies in the literature, we exemplify a case in which a complex phase in the singlet scalar potential is transmitted to the fermion sector via the higher-dimensional operators and drives BAU. We point out that an electric dipole moment of the electron can be suppressed due to the Higgs mass degeneracy and the presence of a new electron Yukawa coupling. Thus, viable parameter space for electroweak baryogenesis is still wide open for the latest experimental bound set by the JILA Collaboration.
I. INTRODUCTION
In the standard model (SM), an explanation of a baryon asymmetry of the Universe (BAU) via electroweak baryogenesis (EWBG) mechanism [1,2] is excluded due to a lack of a strong first-order electroweak phase transition (EWPT) [3] and insufficient CP violation [4].Despite its failure, the mechanism is still attractive from the viewpoint of testability, and the EWBG possibility in various models has been actively investigated in light of experiments, such as the Large Hadron Collider and electric dipole moment (EDM) of the electron (d e ).One viable scenario compatible with the current LHC data is the so-called degenerate scalar scenario in which new scalar masses are close to 125 GeV.Such a scenario can be realized in the SM with a complex scalar (CxSM) and comprehensively studied in connection with dark matter (DM) physics, where it is shown that a spin-independent DM cross section with nucleons is suppressed thanks to the Higgs mass degeneracy [5].Moreover, the scenario can accommodate the strong firstorder EWPT though the suppression mechanism for the DM cross section turns out to be another kind [6].
In the CxSM, even though complex phases can, in principle, exist in the scalar potential, they cannot be the sources for BAU since the SU(2) singlet scalar field does not couple to SM fermions directly.The simplest way to get around this problem is to introduce higherdimensional Yukawa interactions containing the singlet scalar field.If the coefficients of the operators are complex, pseudoscalar interactions would be induced, driving EWBG [7][8][9][10].If the coefficients happen to be real, on the other hand, the CP violation relevant to EWBG should result from the complex phase of the scalar po-tential.Ref. [11] shows that the strength of the firstorder EWPT could be weakened by the complex phase of the scalar potential, but one could still have the strong first-order EWPT compatible with EWBG.On the other hand, the BAU estimate is not conducted there and is left for future work.
Experimental searches for CP violation beyond the SM are essential for probing the EWBG possibility.Currently, the electron EDM is the most sensitive to CP violation.In 2018, ACME Collaboration placed an upper bound on d e as |d ACME e | < 1.1 × 10 −29 e cm at 90% C.L. [12] and in 2022, JILA Collaboration further improved the bound as |d JILA e | < 4.1 × 10 −30 e cm at 90% C.L. [13].While maintaining the BAU, some suppression mechanisms should be present to avoid such an unprecedentedly tight EDM bound (for cancellation mechanisms, see, e.g., Refs.[14,15]).
In this letter, we investigate the EWBG feasibility in the CxSM with higher-dimensional operators.In particular, we consider cases where the complex phase exists in the scalar potential, which is transmitted to the SM fermion sector via the dimension-5 Yukawa interactions with and without complex coefficients.Our study shows that the complex phase of the scalar potential yields the right ballpark value for BAU without resorting to the complex coefficients of the higher-dimensional operators.On the other hand, d e can be suppressed by the presence of the Higgs mass degeneracy and new electron Yukawa coupling, thus evading the latest upper bound from the JILA experiment.
II. MODEL
The CxSM is the extension of the SM by adding a complex SU(2) singlet scalar field (S) [16].In the most general scalar potential, there are 5 real parameters and 8 complex parameters.As a first step toward the general analysis, we take a principle of minimality to simplify our arXiv:2309.09430v2[hep-ph] 22 Nov 2023 analysis, which is also employed in our previous work [6,11]. 1 The scalar potential we consider in this work is given by where Without the a 1 and b 1 terms, V 0 has a global U(1) symmetry, and a massless Nambu-Goldstone boson would appear if the symmetry is spontaneously broken.Moreover, a 1 is necessary to break a Z 2 symmetry S → −S, which dodges a domain wall problem.If the scalar sector preserves CP, V 0 is invariant under the transformation χ → −χ, and χ could be DM.However, as investigated in Ref. [6], the DM relic abundance is too small in the parameter space where EWPT is strong first order.We, therefore, degrade χ to an ordinary unstable particle by allowing CP violation as needed for EWBG.While both a 1 and b 1 can be complex, their relative phase is only physical.As our convention, only a 1 is treated as the complex parameter and parametrized as a 1 = a r 1 + ia i 1 .In this setup, the tadpole (minimization) conditions for h, s, and χ are respectively given by where the symbol ⟨• • • ⟩ denotes that the fluctuation fields are taken zero after their derivatives, and |v S | 2 = v r2 S + v i2 S .After imposing the tadpole conditions, the mass matrix in the basis (h, s, χ) is cast into the form which is diagonalized by an orthogonal matrix O as ).We parametrize the matrix O as where s i = sin α i and c i = cos α i (i = 1, 2, 3).In our work, the 8 original parameters using the tadpole conditions together with the mass condition.We note that a i 1 is given by which implies that a i 1 = 0 if v i S = 0, i.e., the explicit CP violation must be associated with the spontaneous CP violation, but not vice versa.
The Higgs couplings to fermions (f ) and gauge bosons (V = Z, W ± ) are defined as where κ if = O 1i and κ iV = O 1i .Note that the presence of the complex parameters in the scalar potential does not give rise to pseudoscalar coupling in L hi f f , meaning that EWBG is not driven in this setup.To circumvent this issue, we introduce dimensional-5 operators.The relevant terms in the following discussion are where q L denotes the up-type left-handed quark doublet of the third generation, while ℓ L is the down-type lefthanded lepton doublet of the first generation.t R and e R are the right-handed top and electron, respectively.Λ is a cutoff scale and H = iτ 2 H * with τ 2 representing the second Pauli matrix.y t and y e are the top and electron Yukawa couplings in the SM, respectively, while c t and c e are general complex parameters.For later use, we As shown below, c e could be pivotal in suppressing the electron EDM.
Let us redefine the Higgs couplings to the fermions in the presence of the dimension-5 operators as where As seen, the pseudocouplings g P hi f f exist because of the dimension-5 operators, and χ is now interpreted as the pseudoscalar.
Our primary interest is the case in which the complex phase in the scalar potential is the only source for the CP violation that drives EWBG.Secondarily, to what extent complex c t and c e can change the former result.In what follows, we consider the 2 cases: 1.Both c t and c e are real 2.Both c t and c e are complex We make a comment on a case in which c t is complex while c e is real at the end of Sec.V.
Before closing this section, we briefly describe the degenerate scalar scenario that can mimic the SM.For illustration, we consider a process gg in our benchmark points, where Γ hi are the total decay width of h i , we can use a narrow decay width approximation [17,18].With the approximation, the cross section normalized by the SM value is cast into the form where we have used Γ hi ≃ κ 2 iV Γ SM h with Γ SM h being the total decay width of the SM Higgs boson.For |c t | = y t and Λ = 1.0 TeV, the deviation from the SM value would be about 6%, which is still consistent with the current LHC data [19,20]. 2While somewhat lower Λ could be allowed experimentally, detailed collider analysis would be required for that, and we do not pursue this possibility in the current work.We have confirmed that our conclusion does not change even when Λ = 0.5 TeV.
Currently, experimental constraints on the Higgs total decay width are Γ exp h < 14.4 MeV (ATLAS [21]) and Γ exp h = 3.2 +2.4 −1.7 MeV (CMS [22]), which are not precise enough to provide a valuable constraint to our scenario.
2 Note that deviations of other processes such as the Higgs decay to diphoton are also O( ) ∼ 6% in our study, which is consistent with the current LHC data [19,20].
III. ELECTROWEAK BARYOGENESIS
We are following closely the work of Refs.[23][24][25], derive the semiclassical force in the presence of the CP violation discussed in the previous section.The Yukawa interaction with a spacetime-dependent complex mass is defined as where ∂ / = γ µ ∂ µ .Since the thickness of the bubble wall is much smaller than that of the radius, we can approximate it as a planner.In this case, the spacetime dependence of m f is only z which is the coordinate of the perpendicular to the wall.
From the above Yukawa Lagrangian, the equation of motion is given by where . The semiclassical force is found to be where The upper and lower signs correspond to particles and antiparticles, respectively.We also note that particles with opposite spin receive the opposite CP-violating force.
The nonzero momenta parallel to the wall can enhance the CP-violating part, as referred to by Ref. [24].In our case, the top mass during EWPT has the form where ρ(z), ρ r S (z), and ρ i S (z) are the bubble wall profiles parametrized as ⟨H⟩ T = (0 ρ(z)/ √ 2), ⟨S⟩ = (ρ r S (z) + iρ i S (z))/ √ 2, while the phase θ t (z) is expressed as The detail of the bubble wall calculations is given in Ref. [11].After solving transport equations, one can find the baryon-to-photon ratio (η B ) as [25] In our model, dominant corrections to d e come from the so-called Barr-Zee diagrams [27].We decompose them into two parts ) W with the subscripts of the parentheses representing the particle running in the upper loop in the Barr-Zee diagrams, as depicted in Fig. 1.
The top-loop contributions to d e in the degenerate mass limit becomes where τ th = m 2 t /m 2 h with m h ≡ m h1 = m h2 = m h3 , and c t,e = |c t,e |e iϕt,e .f (τ th ) and g(τ th ), are the loop functions defined in Refs.[27].In our convention, e represents the positron charge.Eq. ( 29) implies that (d hγ e ) t vanishes when ϕ t = ϕ e + nπ with n being the integer, let alone c t and c e are both real.
The W -loop contributions to d e are induced by the complex c e , which have the form where J γ W (m hi ) denotes the loop function [28], and one can find where
V. NUMERICAL RESULTS AND DISCUSSIONS
As studied in Ref. [11], 0.3 ≲ v i S ≲ 0.5 is the range where the first-order EWPT is strong enough to suppress baryon-changing processes and bubble nucleation happens.Since the first-order EWPT is driven by a tree-level potential barrier, its strength would remain unchanged even after including the dimension-5 Yukawa operators (12).We take a parameter set BP1 adopted in Ref. [11] for illustrative purposes but with the sign of v i S being flipped.The inputs and outputs are summarized in Table III.Regarding c f (f = t, e), we set |c f | = y f and take ϕ f as the free parameters.
In the case of ϕ t = ϕ e = 0, CP violation solely comes from the scalar potential.With this CP violation, we calculate the BAU in the cases of Λ = 1.0, 1.5, and 2.0 TeV, respectively.The results are summarized in Table II, where |d e | and its details are also shown.One can see that the Λ = 1.0 TeV case yields η B = O(10 −10 ), while the other two cases provide the smaller η B to some extent.Even though the obtained values of η B are somewhat insufficient for explaining the observed one, we make no strong claims about the numbers since the perturbative calculations of EWPT and BAU employed in this work are generally subject to significant theoretical uncertainties.Further theoretical improvements should be left to future work.Now, we discuss the case of complex c t and c e .In this case, there are three sources for CP violation, and v i S and c t are responsible for EWBG.Finally, some comments are noted.
• One may ask whether the cancellation of the electron EDM can occur in concert with the complex c t without resorting to the phase alignment with c e .In principle, this can happen.However, this type of cancellation becomes effective only when the scalar masses are not close to each other.
• Other EDMs such as neutron and Mercury could be significant in exploring this scenario.In doing so, however, it is necessary to introduce additional new Yukawa couplings of the first-generation quarks.This topic should be studied separately from the present analysis.
• Instead of the dimension-5 operators, we could consider dimension-6 Yukawa interactions, such as From the dimensional analysis, CP violation in this case would be more suppressed than in the dimension-5 operator case.It is found that η B < 1.0 × 10 −10 and |d e | < 1.0 × 10 −30 e cm for the same parameter set as in the dimension-5 operator case.In this case, the EDM suppressions due to the additional factor 1/Λ and scalar mass degeneracy are strong enough to avoid the EDM bounds, and the phase alignment ϕ t = ϕ e + nπ is not necessarily required.
• In the general scalar potential, we have more complex parameters coming from S 3 , SH † H, etc.In such an enlarged parameter space, the EDM cancellation would be more effective, while the BAU may be more enhanced.
• Double Higgs production processes are one of the interesting collider signatures of EWBG.As mentioned in Sec.II, the modification by the top Yukawa couplings is typically 6%.On the other hand, the triple Higgs couplings in this model could get large compared to the SM value.Among all the triple Higgs couplings λ hihj h k (i = 1, 2, 3), we find that λ h1h1h1 is the largest in our benchmark point, which is about 1.4 times larger than that in the SM.Even though the current LHC cannot measure the triple Higgs coupling [20,29], future colliders may be capable.We defer the detailed analysis to future work.
VI. CONCLUSION
We have studied the possibility of EWBG in the CxSM with the dimension-5 Yukawa interactions.We consider the two cases: one is the case in which CP violation arises only from the scalar potential and propagates to the SM fermion sector by the dimension-5 top Yukawa interaction, and the other is the case where the coefficient of the dimension-5 Yukawa interaction additionally yields CP violation.It is found that the former leads to η B = O(10 −10 ), and the additional CP violation in the latter helps to increase η B to some extent.Even though the nominal values of η B in our benchmark points are smaller than the observed value by a factor of a few, the deficit might be compensated by theoretical uncertainties that could reside in the perturbative treatments of EWPT and BAU.A more elaborate analysis will be left to future research.
We also investigated the electron EDM in the two cases mentioned above.The electron EDM is suppressed due to the Higgs mass degeneracy, and the ACME and JILA constraints can be evaded for the real c t and c e cases.In contrast, in the complex c t and c e case, the phase alignment ϕ t = ϕ e + nπ is additionally needed to be consistent with the experimental bounds.
In conclusion, the EWBG parameter space in our scenario is still wide open after the recent EDM updates.
FIG. 1 .| < 4 . 1 ×
FIG. 1. Dominant two-loop contributions to de.The left diagrams are denoted as (d hγ e )t and (d hZ e )t, while the right ones as (d hγ e )W and (d hZ e )W .
FIG. 2 .
FIG. 2. The electron EDM as a function ofh 2 in the case that |ct| = yt, |ce| = ye, ϕt = ϕe = 0, and Λ = 1.0 TeV.We take the parameter set given in Table III while m h 2 is treated as the free parameter.Here, d t e and d W e are the twoloop contributions to the electron EDM, depicted as the left and right diagrams in Fig. 1, respectively.
example, we take Λ = 1.0 TeV, |c t | = y t , |c e | = y e , ϕ t = ϕ e = 0, and thus CP violation arises from the nonzero v i S .Fig. 2 shows |d e | (green solid line) and its details |d t e | (blue dotted line) and |d W e | (orange dashed line) against m h2 .The upper dotted horizontal line denotes the experimental bound of ACME, while the lower one represents the JILA bound.As discussed above, the both |d t e | and |d W e | would be suppressed as m h2 approaches 125 GeV(= m h1 ), evading ACME and JILA constraints.This example clearly illustrates that the degenerate scalar scenario simultaneously provides an exquisite parameter space compatible with the LHC and the electron EDM data.
Fig. 3 displays η B and |d e | in the (ϕ t , ϕ e ) plane.The vertical dotted lines denotes
Therefore, regardless of c e , (d hγ e ) W vanishes when m h1 = m h2 = m h3 .Similarly, we can obtain the same vanishing conditions for (d hZ e ) t and (d hZ e ) W .The conditions for the vanishing d t e and d W e are summarized in Table I.To see the suppression behavior numerically, a typical example is given here.The input parameters are summarized in Table III but with m h2 being free.In this ηB/10 −10 |de|/10 −30 d t e /10 −30 d W e /10 −30
TABLE II .
Summary of ηB and |de| in the case of |ct| = yt, |ce| = ye, and ϕt = ϕe = 0.The electron EDM is given in units of e cm. | 4,559.8 | 2023-09-18T00:00:00.000 | [
"Physics"
] |
Comparison of the Efficacy of Different Techniques for the Removal of Root Canal Filling Material in Artificial Teeth: A Micro-Computed Tomography Study
This study aimed to assess the efficacy of canal filling material removal using three different techniques after filling with a Gutta–Percha (GP) cone and calcium silicate-based sealer, by measuring the percentage of volume debris of GP and sealer remaining intracanal with micro computed tomography (micro-CT). The filling material was removed from 30 plastic teeth by a nickel–titanium (Ni–Ti) rotary retreatment system. Final irrigation was performed with 2 mL of saline and 10 specimens were randomly allocated to a conventional group. In the passive ultrasonic irrigation (PUI) group, ultrasonic irrigation was added to the conventional group (n = 10). In the Gentlefile Brush (GF Brush) group, irrigation with GF Brush was added to the conventional group (n = 10). Remaining filling material was measured using micro-CT imaging analysis. The total mean volume of residual filling material after retreatment in the conventional group, PUI group and GF Brush group were 4.84896 mm3, 0.80702 mm3, and 0.05248 mm3, respectively. The percentage of filling material remaining intracanal was 6.76% in the conventional group, 1.12% in the PUI group and 0.07% in the GF Brush group. This study shows that the cleaning effect of the GF Brush system is superior to those of Ni–Ti retreatment files and the PUI system in the apical area.
Introduction
The goal of endodontic treatment is the eradication of harmful microorganisms from the root canal. Thus, cleaning and shaping are key for the success of endodontic treatment. However, the anatomical complexities of the root canal system and limitations in current preparation and irrigation techniques lower the success rates for endodontic treatment. Studies concerning the morphology of the root canal system have shown wide variances in the canal shape and the presence of two or more canals in a single root. Furthermore, complete disinfection in the presence of several curvatures and narrow canals is difficult to achieve by all known techniques, whether chemical or mechanical. Consequently, the reported success rate for root canal treatment (RCT) is approximately 75% [1]. Although RCT is a reliable and highly successful treatment, some cases do exhibit post-treatment disease.
Nonsurgical RCT is the first option for the treatment of postendodontic disease. The retreatment procedure is mostly similar to the initial RCT procedure, with the greatest difference being the removal procedure for the root canal filling material during retreatment. There could be some necrotic tissue or bacteria among the filling material, which potentially cause persistent inflammation and pain. Therefore, dental materials in the root canal system should be completely removed in the initial step of retreatment. However, filling materials such as gutta-percha (GP) and sealers are difficult to remove because they are trapped within the irregular root canal system. Several techniques have been used for the removal of filling material from root canals; these include stainless steel (SS) hand files, nickel-titanium (Ni-Ti) rotary instruments, and ultrasonic tips [2,3]. Rotary instruments are widely used and reportedly remove filling material in a safe and efficient manner, with high success rates [4,5]. Nevertheless, studies have shown that none of these retreatment procedures can completely clean the root canal wall, particularly in the apical third [6,7]. Selection of an instrument that can effectively clean the GP and sealer debris in the apical third of root canals is very important, considering most instruments are generally interrupted by the various curves within the canal system. In addition to rotary instruments, ultrasonic instruments are used as auxiliary tools for cleaning root canals. Passive ultrasonic irrigation (PUI) has the potential to remove dentinal debris, organic tissue, and calcium hydroxide from inaccessible root canal areas [8,9]. Grischke et al. reported that an ultrasonic irrigation protocol was superior to other techniques investigated for the removal of sealer from the root canal surface during endodontic retreatment [10]. In other studies, the use of PUI after mechanical instrumentation ensured more efficient material removal during endodontic retreatment than did other techniques such as chloroform irrigation, xylene irrigation, and eucalyptol irrigation [11,12].
Recently, a new SS system known as Gentlefile (GF; MedicNRG, Kibbutz Afikim, Israel) was released. Even though the instruments are made of SS, they have shown better mechanical properties relative to those of Ni-Ti instruments. Moreinos et al. reported that the GF system required longer time and more rotations to fracture compared with the ProTaper and RevoS systems, and the GF system applied less vertical force to the canal in comparison with the ProTaper and RevoS systems [13]. The GF system also offers a brush (GF Brush) comprising six SS strands that automatically open outwards when operated in a handpiece with a speed of 6500 rpm. The original aim of the GF Brush is to aid irrigation, but it is expected to show excellent efficiency for the removal of substances from root canals because of its design. Neelakantan et al. examined the effectiveness of irrigant agitation with the GF Brush after root canal preparation and concluded that the use of the GF Brush resulted in significantly less pulp tissue remnant compared with syringe irrigation [14]. A combination of flexibility and centrifugal movement would facilitate access and cleaning in irregular parts. In case of retreatment, the GF Brush is expected to aid in the removal of GP and sealer particles stuck on the canal walls. Although there is a study on the efficacy of the GF Brush in initial root canal treatment, no studies have compared the removal efficiency of canal filling material using the GF Brush.
The aim of this study was to assess the efficacy of canal filling material removal using three different techniques after filling with a GP cone and calcium silicate-based sealer, by measuring the percentage of volume debris of GP and sealer remaining intracanal with micro computed tomography (micro-CT). The null hypothesis was that the GF Brush system would demonstrate similar efficacy in removal of canal filling material as the retreatment Ni-Ti file and PUI system.
Preparation of Tooth Samples
Thirty-three artificial teeth made of plastic (TrueTooth, Dental Cadre, Santa Barbara, CA, USA) were used for this study (Figure 1a). The teeth were customized samples reproducing the shape of the human mandibular first premolar and exhibiting access openings with a type I canal as per Weine's classification [15]. A #15 K-file (Dentsply Maillefer, Ballaigues, Switzerland) was inserted into the canal to determine the working length (WL), which was recorded as 21 mm. A #40 master apical file was used. The canal in all samples was instrumented using the ProTaper Next Ni-Ti system (Dentsply Maillefer) coupled with the Dentsply X-Smart Plus motor (Dentsply Maillefer). According to the manufacturer's instructions, X1, X2, and X3 files were used up to the full WL. Between instruments, each canal was irrigated using distilled water via a 27-gauge needle (Korean Vaccine Co., Seoul, Korea). After the instrumentation was complete, all canals were dried with #25 paper points (Dentsply Maillefer). Subsequently, three teeth were randomly selected for micro-CT, and the acquired images were overlapped for the confirmation of consistency in the prepared canal space. The other 30 teeth were prepared for obturation.
into the canal to determine the working length (WL), which was recorded as 21 mm. A #40 master apical file was used. The canal in all samples was instrumented using the ProTaper Next Ni-Ti system (Dentsply Maillefer) coupled with the Dentsply X-Smart Plus motor (Dentsply Maillefer). According to the manufacturer's instructions, X1, X2, and X3 files were used up to the full WL. Between instruments, each canal was irrigated using distilled water via a 27-gauge needle (Korean Vaccine Co., Seoul, Korea). After the instrumentation was complete, all canals were dried with #25 paper points (Dentsply Maillefer). Subsequently, three teeth were randomly selected for micro-CT, and the acquired images were overlapped for the confirmation of consistency in the prepared canal space. The other 30 teeth were prepared for obturation. Obturation was performed using the single-cone technique. The canals were first coated with a calcium silicate-based sealer (Well-Root ST sealer, Vericom, Chuncheon-si, Korea) via a 24-gauge needle tip provided by the manufacturer. The tip was slowly pulled toward the orifice from the point of engagement in the canal. Then, medium to large GP cones (DiaDent, Cheongju-si, Korea) were customized to size 40 using a GP gauge (Dentsply Maillefer), and a single cone was inserted in each canal. The cone was gently moved with an up-down motion three times to facilitate better penetration of the sealer into finer structures, following which it was cut at the orifice level using System B (SybroEndo, Orange, CA, USA) and vertically condensed. The access cavities were filled with Caviton (GC Corporation, Tokyo, Japan). All filled samples were stored in a humidified chamber (Changshin Science, Seoul, Korea) at 100% relative humidity and 37 °C for 7 days until retreatment. All procedures were performed by a single operator.
Root Canal Retreatment
The temporary filling was removed with a round bur, and the specimens were randomly allocated to three different groups (n = 10 per group) according to the material removal technique.
Conventional Group
The root canal filling material was removed using ProTaper Universal retreatment files (Dentsply Maillefer) according to the manufacturer's instructions. D1, D2, and D3 files were sequentially used with the crown-down technique until WL was reached. The files were manipulated with a brushing action at a constant speed of 500 rpm, as recommended. A solvent was not used. Retreatment was considered complete when no GP/sealer was visible on the surface of the instruments. Root canal refinement was accomplished using #25, 30, 35 and 40 K-files up to WL. Between instruments, the canal was irrigated with 2 ml of distilled water via a 27-gauge needle Obturation was performed using the single-cone technique. The canals were first coated with a calcium silicate-based sealer (Well-Root ST sealer, Vericom, Chuncheon-si, Korea) via a 24-gauge needle tip provided by the manufacturer. The tip was slowly pulled toward the orifice from the point of engagement in the canal. Then, medium to large GP cones (DiaDent, Cheongju-si, Korea) were customized to size 40 using a GP gauge (Dentsply Maillefer), and a single cone was inserted in each canal. The cone was gently moved with an up-down motion three times to facilitate better penetration of the sealer into finer structures, following which it was cut at the orifice level using System B (SybroEndo, Orange, CA, USA) and vertically condensed. The access cavities were filled with Caviton (GC Corporation, Tokyo, Japan). All filled samples were stored in a humidified chamber (Changshin Science, Seoul, Korea) at 100% relative humidity and 37 • C for 7 days until retreatment. All procedures were performed by a single operator.
Root Canal Retreatment
The temporary filling was removed with a round bur, and the specimens were randomly allocated to three different groups (n = 10 per group) according to the material removal technique.
Conventional Group
The root canal filling material was removed using ProTaper Universal retreatment files (Dentsply Maillefer) according to the manufacturer's instructions. D1, D2, and D3 files were sequentially used with the crown-down technique until WL was reached. The files were manipulated with a brushing action at a constant speed of 500 rpm, as recommended. A solvent was not used. Retreatment was considered complete when no GP/sealer was visible on the surface of the instruments. Root canal refinement was accomplished using #25, 30, 35 and 40 K-files up to WL. Between instruments, the canal was irrigated with 2 mL of distilled water via a 27-gauge needle (Korean Vaccine Co.). Final irrigation was performed with 5 mL of distilled water for 30 s. Finally, the root canal was dried with paper points and stored in a dry environment for micro-CT analysis.
PUI Group
In the PUI group, ultrasonic irrigation was added to the procedure described for the conventional group. Ultrasonic irrigation was performed using an ultrasonic endodontic tip (Endosonic Blue, Maruchi, Chuncheon-si, Korea; Figure 1b), which was inserted into the root canal up to 1 mm short of WL and oscillated toward the apex [16]. Activation with 3 ml of distilled water for 60 s was performed three times (total 3 min per tooth). The distilled water was replenished between each activation cycle.
GF Brush Group
In the GF Brush group, irrigation with the GF Brush was added to the procedure described for the conventional group. Final instrumentation was performed using the GF Brush (Figure 1c,d), which was inserted into the root canal up to 1 mm short of WL. Activation with 3 mL of distilled water for 60 s was performed three times (total 3 min per tooth). The distilled water was replenished between each activation cycle.
Micro-CT Analysis and Stereomicroscopy
Micro-CT and image analysis were performed as previously described [17]. A high-resolution micro-CT scanner (SkyScan 1173, Bruker, Billerica, MA, USA) was used to scan the samples in the three groups. All acquired images were reconstructed using NRecon software, version 1.6.6.0 (Bruker microCT, Kontich, Belgium). For evaluation of the residual filling material, three-dimensional (3D) images of the filling material were visualized by surface-CT-Vol (SkyScan). The CT-An software (SkyScan) was used to measure the volume of the prepared canal space in the three sample specimens and the volume of the residual filling material after retreatment in the remaining 30 specimens. The apical region was defined as the area between 1 and 5 mm from the apex, the middle region as the area between 5 and 10 mm from the apex, and the coronal region as the area between 10 and 15 mm from the apex.
After micro-CT, the teeth were longitudinally sectioned at the labial and lingual surfaces using a low-speed diamond wheel (Struers Minitom, DK-2610, Rodovre, Denmark) under water cooling. Each root surface was observed under a stereomicroscope (Zeiss, Gottingen, Germany).
Statistical Analysis
The amount of residual filling material was expressed as a percentage of the total area of each section in the root canal. The measurements were evaluated by a single observer blinded to the study groups. The Shapiro-Wilk test was used to verify whether the data were normally distributed. The Student's t test was used to compare the percentage volume of residual filling material in the apical, middle, and coronal regions. A p-value of <0.05 was considered statistically significant. All statistical analyses were performed using SPSS software, version 23 (SPSS, Chicago, IL, USA).
Results
The volume of the prepared canal space was similar in the three sample specimens (71.68401 mm 3 ). The total mean volume of residual filling material after retreatment in the conventional, PUI, and GF Brush groups was 4.84896, 0.80702, and 0.05248 mm 3 , respectively, with percentage values of 6.76%, 1.13%, and 0.07%, respectively (Table 1; Figure 2). In the conventional group, the filling material was distributed evenly on the root canal walls. In the PUI group, the filling material debris was mostly concentrated in the apical region, which was beyond the canal curvature. In the GF Brush group, all three regions showed a small amount of residual material. Stereomicroscopic observation demonstrated that residual filling material stayed on the root canal walls. In the conventional group, the material was uniformly distributed in the three regions, while the PUI and GF Brush groups exhibited an insignificant amount of residual material. In the PUI group, most of the residual material was concentrated in the apical region ( Figure 3). Stereomicroscopic observation demonstrated that residual filling material stayed on the root canal walls. In the conventional group, the material was uniformly distributed in the three regions, while the PUI and GF Brush groups exhibited an insignificant amount of residual material. In the PUI group, most of the residual material was concentrated in the apical region ( Figure 3).
Discussion
In the present study, we compared the efficacy of filling material removal during retreatment in artificial mandibular premolars between the GF Brush, conventional Ni-Ti, and PUI systems. The reproduction of the clinical situation may be regarded as the major advantage of the use of natural teeth for experiment. However, the wide range of variations in three-dimensional root canal morphology makes standardization difficult between groups. On the other hand, the use of artificial teeth allows standardization of degree, location, and radius of root canal curvature in threedimensions [18]. Thus, artificial teeth were selected to obtain more reliable results in terms of the volume of residual filling material after retreatment. Although uniform samples were used, the prepared canal shapes may have differed. Therefore, to determine consistency, we randomly selected three prepared teeth and subjected them to micro-CT. The acquired images were superimposed, and the results showed negligible variations in the prepared canal space among the teeth. This was probably because the canal space in artificial teeth is quite wide, even before canal preparation. Larger files were not used because we were not focused on comparing the preparation efficiency of the instruments. Consequently, we could maintain a consistent canal space volume in all specimens.
We found that the amount of residual filling material was smaller in the PUI group than in the conventional group, with the exception of the apical region. This result was similar to that of a previous study, where the PUI technique was found to be more effective than the conventional technique for the removal of root canal filling material from the cervical and middle thirds during endodontic retreatment [19]. Additional PUI retreatment eliminates sealer and GP debris beyond the curvature or in areas that cannot be accessed by conventional retreatment files. During PUI, free intracanal movement of the file is necessary for easy penetration of the solution into the root canal system and a more powerful cleaning effect [20]. During PUI, energy is transmitted from a file or smooth oscillating wire to the irrigant by means of ultrasonic waves that induce two physical phenomena: stream and cavitation of the irrigant. The acoustic stream can be defined as rapid movement of the fluid in a circular or vortex shape around the vibrating file, while cavitation is defined as the generation of steam bubbles or the expansion, contraction, and/or distortion of preexisting bubbles in a liquid [21]. The effect of an ultrasound tip is the maximum in the middle and
Discussion
In the present study, we compared the efficacy of filling material removal during retreatment in artificial mandibular premolars between the GF Brush, conventional Ni-Ti, and PUI systems. The reproduction of the clinical situation may be regarded as the major advantage of the use of natural teeth for experiment. However, the wide range of variations in three-dimensional root canal morphology makes standardization difficult between groups. On the other hand, the use of artificial teeth allows standardization of degree, location, and radius of root canal curvature in three-dimensions [18]. Thus, artificial teeth were selected to obtain more reliable results in terms of the volume of residual filling material after retreatment. Although uniform samples were used, the prepared canal shapes may have differed. Therefore, to determine consistency, we randomly selected three prepared teeth and subjected them to micro-CT. The acquired images were superimposed, and the results showed negligible variations in the prepared canal space among the teeth. This was probably because the canal space in artificial teeth is quite wide, even before canal preparation. Larger files were not used because we were not focused on comparing the preparation efficiency of the instruments. Consequently, we could maintain a consistent canal space volume in all specimens.
We found that the amount of residual filling material was smaller in the PUI group than in the conventional group, with the exception of the apical region. This result was similar to that of a previous study, where the PUI technique was found to be more effective than the conventional technique for the removal of root canal filling material from the cervical and middle thirds during endodontic retreatment [19]. Additional PUI retreatment eliminates sealer and GP debris beyond the curvature or in areas that cannot be accessed by conventional retreatment files. During PUI, free intracanal movement of the file is necessary for easy penetration of the solution into the root canal system and a more powerful cleaning effect [20]. During PUI, energy is transmitted from a file or smooth oscillating wire to the irrigant by means of ultrasonic waves that induce two physical phenomena: stream and cavitation of the irrigant. The acoustic stream can be defined as rapid movement of the fluid in a circular or vortex shape around the vibrating file, while cavitation is defined as the generation of steam bubbles or the expansion, contraction, and/or distortion of pre-existing bubbles in a liquid [21]. The effect of an ultrasound tip is the maximum in the middle and coronal regions because of the direction of operation [8]. In our samples, the apical region was located beyond curvature from the canal orifice. This may have caused limited movement of the tip during its passage through the curvature, resulting in reduced cleaning efficiency in the apical third.
Compared with the conventional and PUI groups, the GF Brush group showed a significantly smaller volume of residual filling material in the apical region. Almost 99.76% of the filling material in the apical region was eliminated by this system. The superior effects of the GF Brush system may be attributed to the mechanism of the GF Brush. When the brush is not rotating, the strands are present in a twisted form. However, at a high rotating speed, the strands open to cover the entire canal diameter. These opened strands mechanically remove the remaining debris and bring the irrigant into intimate contact with the canal surface. Flexibility and centrifugal movement facilitate access and cleaning in irregular parts. These characteristics of the GF Brush system may contribute to the better cleaning efficiency throughout the canal system, particularly the apical region. To achieve more effective cleaning and shaping of the root canal system, the endodontist must ensure cleanliness of the apical third of the canal, and we believe that the GF Brush system is an effective tool to achieve this goal.
This study has some limitations. First, the artificial teeth used in this study cannot perfectly replicate natural teeth due to complex root canal anatomy such as lateral canal, isthmus, and fin. In addition, the microscopic structures of the dentine of natural teeth are absent in the artificial teeth; therefore, the adhesion between the endodontic filling material and the root canal wall cannot be reproduced.
Within the limitation of this in vitro study, our findings suggest that the GF Brush system is superior to conventional Ni-Ti retreatment files and PUI in terms of effective removal of root canal filling material during retreatment, particularly from the apical third of the canal. This is because the cleaning effect of the GF Brush is not impeded by the curvature of the canal because of its design. Additional experiments in various shapes of root canal and/or natural teeth may be helpful in clinical applications. | 5,317.2 | 2019-07-01T00:00:00.000 | [
"Medicine",
"Materials Science"
] |
EXPERIMENTAL AND NUMERICAL STUDIES ON WAVE TRANSFORMATION OVER ARTIFICIAL REEFS
A laboratory measurement on the flow field, turbulence and wave energy of spilling breakers over artificial reefs is presented. Instantaneous velocity fields of propagating breaking waves on artificial reefs were measured using Particle Image Velocimeter (PIV) and Bubble Image Velocimeter (BIV). Variations of water surface elevation were observed by using Charge Coupled Device (CCD) cameras with horizontal posture. The experimental results showed that the initial bubble velocity in the aerated region is faster than phase speed with a factor of 1.26. The velocity profiles are identical to the shallow water theory. It is found that a low flow velocity exists due to an opposite but equal onshore and offshore velocity. Significant turbulent kinetic energy and turbulent Reynolds stress are produced by breaking waves in the front of aerated region, then move offshore and decay. The calculated total energy dissipation rate was compared to that based on a bore approximation. 2
INTRODUCTION
Hard engineering coastal structures such as coastal dyke, offshore breakwaters and groin are built to protect the coastal line. However, these structures have disadvantages, such as coastal erosion by improper design that accelerate disappearance of sand, decrease of attraction and harmonious with environment. Along with the rising of environmental consciousness in recent years, the hard engineering methods are not only solutions to coastal defense. The flexible working methods are now becoming alternate solutions such as submerged breakwater, artificial nourishment and artificial reefs. Among them, artificial reefs have high potential in practical applications because they are acting like natural reefs. In the past, many researchers investigated the interaction between wave and structures through numerical or experimental test. In the last decade, non-intrusive measurement techniques were employed to observe the various complex flow phenomena under breaking waves over coastal structures. The laser Doppler velocimetry (LDV) and particle image velocimetry (PIV) technique were successfully used to measure breaking wave flow field in a surf zone. Ting and Kim (1994) applied the LDV to observe the vortex generation in water waves propagation over a submerged obstacle, and compared the scale of vortex by K-C number.They found that the scale of the vortex above the structure surface was affected by the K-C number. Petti et al. (1994) investigated the wave velocity field measurement over a submerged breakwater by PIV method, and discussed the variations of vortex under the wave breaking condition. Chang el al.(2001) used the experiment and numerical simulation to investigated vortex generation and evolution in water waves propagating over a submerged rectangular obstacle, and measured the flow field around the submerged obstacle by PIV method.
In wave breaking condition, Yasuda el at. (1997) investigated the flow field and their breaker types by simulating the solitary wave propagation over a submerged obstacle. Jansen (1986) used the fluorescent dye and ultraviolet light to measure displacement of the particles in aerated region. However, the results suffered from poor resolution for spatial variation. Chang and Liu (1998) used the PIV to measure the maximum velocity and associated acceleration and vorticity of the overturning jet of a breaking wave. Unfortunately,PIV is limited in the case where breaking wave produced the bubbles and thus scattered the laser light. PIV is only valid in the water region, therefore the instruments were limited by the non-aerated region. Greated and Emarat (2000), Kirby (1994, 1995), Perlin et al. (1996) used the LDV (Laser Doppler Velocimetry) to measure the aerated region, but it is only for single-point measurement.
To overcome air bubble effect on the measurement Hassan et al. (1998), Nishino et al. (2002 and Lindken and Merzkirch (2002), applied the 'shadowgraphy' method to to measure bubble velocity by correlating bubbles or tracking each bubble in the recorded images. Govender et al., (2002) and Ryu et al. (2005) also have successfully measured the velocity field in the aerated region and in the overtopping region with air bubble entrainment. In this study, the mechanism of two-dimensional flow field of spilling breakers on an artificial reef was investigated experimentally using bubble image velocimetry. The experiments were conducted in a wave flume at Department of Hydraulic and Ocean Engineering of National Cheng Kung University. The wave tank is 25 m long, 0.5 m wide and 0.6m high. The wave maker is of piston type installed at one end of the wave tank and controlled by a computer. The wave absorber is at the other end of the tank to absorb the wave energy and reduce reflection. The water depth was kept constant at 0.165 h = m high. A rectangular model structure is 2.4 m long and 0.115 m high. The wave data were recorded at a sampling rate of 100 Hz by capacitancetype wave gauges distributed at 4 fixed locations. Fig.1 is a schematic diagram of the facilities and apparatus layouts.
EXPERIMENTAL FACILITIES AND SETUP
The PIV system used in the present study includes a dual-head pulsed laser, laser light sheet optics, a CCD camera, and a synchronizer. The dual-head pulsed laser is Nd: YAG laser that has a 20 Hz repetition rate and 120 mJ/pulse maximum energy output. It was used as the PIV illumination source. Images were recorded using a 12-bit CCD camera that has a 1600 × 1200 pixel resolution and 30 frames per second (fps) maximum framing rate. The BIV system used in the present study includes a high speed camera, and two 600W light bulbs. The images were captured by IDT MotionProX3TM PLUS high speed camera. The camera has a resolution of 1280 × 1024 pixels and a maximum framing rate of 2000 fps.
The water in the wave flume was seeded with nearly neutrally buoyant hollow glass spheres particles (TSI, normal mean diameter: 8-12 m µ ; density: 1.05-1.15 3 g cm ) to enhance illumination efficiency during the PIV measurements. Because the wave-breaking region in the surf zone is too large to be captured in one single frame, the complete spatial distribution of velocities was integrated by the mosaic of frame from 5 fields of view (FOV), as shown in Fig2. Note that the origin x = 0 is at the structure front wall, and the size of each FOV is shown in Table 1 and Table 2.
VALIDATION OF THE BIV METHOD
Because there was no particular BIV measurement system software package, the validation used the open source software created by Nobuhito Mori. Originally MPIV program was designed to analyze PIV. With modification of Mori's MPIV program,the software can be applied to analyze BIV as well. In the present study, the validation of the accuracy of this software is performed by the captured bubble image as shown in Fig. 3. After sampling 10 sets of images, the velocity calculating by modified MPIV is and from the images directly. The velocity was calculated using modified MPIV referred as U, and from the images directly is referred as U R . Comparison between each velocity measurement was made to validate the BIV techniques. The result of the validation of the BIV method was shown in Fig. 4. The horizontal axis represents the time interval images; the vertical axis indicates the actual value of the error percentage. The mean error is about 2.10% and the maximum error is about 3.89%. It proved the accuracy of the bubble velocity measurement.
RESULTS AND DISCUSSION
In the present study, the incident wave period is 1 s and wave height is 4 cm. There were four wave gauges located at 8m, 0m, -0.75m and -1.75m from the front wall of the test model. Considering the test region and the resolution of the CCD camera, the test region was divided into five FOVs. Fig.5 shows that the surface elevation measured at x = 8 m from the front wall of the test model. The figure shows the wave become steady at the eighth wave, so the images are captured from the eighth to eleventh waves. The waves are not affected by the reflection of the slope at the end of the water tank and the secondary reflection from the panel of the wavemaker. Fig. 6 shows the result of combined flow field of five regions as mentioned above. At t = 1/10T the breaking waves pass the tip of the artificial reef, the velocity near the free surface is larger than the velocity near the structure. It is shown that the flow field in the upper layer is slightly affected by the velocity of wave trough in offshore direction. Because of the offshore velocity, the flow field in the lower layer forms a convergent stagnation point, and the velocity approaching to zero. At t = 4/10 T, the wave form reached its stability value and started to break. From t =5/10T to t =10/10T, the waves are broken and the water body affected by the gravity trapping air inside it which also causes resistance to the flow. In the aerated region near the structure of the lower layer, stagnation area occurs which obviously different from the upper layer near the free surface. As time increased, the stratification phenomena area become more obvious and the stagnation area becomes larger as losing its momentum caused by bubble collision during the wave breaking.
In Fig. 7 the color map stands for the u/C value, u means horizontal flow velocity and C means theoretical wave velocity. When u/C value is larger than 1, the wave starts to break. At t = 4/10T, the u/C value is larger than 1, so the wave starts to break. The maximum value of u/C can reach about 2 at t = 5/10T, and most of the values can reach 1.4~1.8. At t = 6/10T to t = 10/10T, the u/C value decrease rapidly in the aerated area. After wave breaks, the wave trap air into the water and affected by the front and bottom water body. The velocity starts to decrease, and the u/C value closes to 1. The numerical result calculated by Flow 3D was shown in Fig. 8. At t = 7/10T the stagnation area was occurred and move with wave flow. At t = 2/10T to t = 6/10T the stratification phenomena become more obvious, the numerical results show good agreement with experimental results.
CONCLUSION
In this study, bubble image velocimetry and particle image velocimetry combining with highspeed CCD cameras and other devices are usedto record the interaction between wave breaking and artificial submerged reef model. The wave transformations over artificial reef are investigated. The wave elevation measured by the wave gauge above the artificial reef shows the wave height decay obviously when wave breaks above the artificial reef. After the wave breaking, there was a stratification between the water body near the surface and the structure, and the latter has a smaller velocity. As time increased, the bubble range became narrow and the stratification grew into more obvious. When the wave is breaking, the bubble has the maximum velocity. The velocity within the aerated area is affected by the front and bottom water body, the u/C value form 2 decrease to 1.2~1.4, then declines to about 1. Both experimental and numerical results showed the wave height decayed after the wave breaking and it became more obvious as the distance further from the wave breaking. | 2,612.4 | 2011-01-25T00:00:00.000 | [
"Physics",
"Engineering"
] |
Gold-Catalyzed Synthetic Strategies towards Four-Carbon Ring Systems
: Four carbon ring systems are frequently present in natural products with remarkable biological activities such as terpenoids, alkaloids, and steroids. The development of new strategies for the assembly of these structures in a rapid and e ffi cient manner has attracted the interest of synthetic chemists for a long time. The current research is focused mainly on the development of synthetic methods that can be performed under mild reaction conditions with a high tolerance to functional groups. In recent years, gold complexes have turned into excellent candidates for this aim, owing to their high reactivity, and are thus capable of promoting a wide range of transformations under mild conditions. Their remarkable e ffi ciency has been thoroughly demonstrated in the synthesis of complex organic molecules from simple starting materials. This review summarizes the main synthetic strategies described for gold-catalyzed four-carbon ring formation, as well as their application in the synthesis of natural products. tandem rearrangement/cyclization indolyl 66.
Unsaturated systems activated by gold(I) neutral and cationic complexes can trigger a wide range of nucleophilic additions, processes often understood as cationic stepwise cascade mechanisms in which vinyl-gold species are commonly generated ( Figure 2). These intermediates can further react with other nucleophiles, such as double or triple bonds, oxygenated functional groups, or strained cycles through 1,2-alkyl rearrangements. In this context, gold(I) complexes have emerged as efficient catalysts towards four-carbon rings synthesis. As of now, mainly two approaches to access these carbocycles using gold(I) catalysts have been developed: [2+2] cycloadditions and ring expansions. Interestingly, the key step of biosynthetic routes oriented to obtain these rings relies on carbocationic cyclization/cycloisomerization processes, which in general involve [2+2] cycloadditions or occasionally ring expansions, either in a concerted or a stepwise fashion [100].
In the same way, phosphoramidite ligands have been shown to be useful for these diastereoand enantioselective gold(I)-catalyzed cycloadditions of allenenes (Scheme 2) [118][119][120]. A computational study supports a stepwise mechanism in which the alkene adds to the gold coordinated allene 6, generating a vinyl-gold intermediate. Two possible stereochemical rearrangements through cis-7 and trans-9 can be proposed considering the relative position of substituents on the formed cyclopentane ring. In both pathways, an interaction between the In the same way, phosphoramidite ligands have been shown to be useful for these diastereo-and enantioselective gold(I)-catalyzed cycloadditions of allenenes (Scheme 2) [118][119][120]. A computational study supports a stepwise mechanism in which the alkene adds to the gold coordinated allene 6, generating a vinyl-gold intermediate. Two possible stereochemical rearrangements through cis-7 and trans-9 can be proposed considering the relative position of substituents on the formed cyclopentane ring. In both pathways, an interaction between the carbocation and gold forming a five membered metallacycle is suggested by the calculations. On the cis pathway, the intermediate evolves to the formation of cyclobutene derivatives 8 through demetalation, whereas carbocationic intermediate trans-9 is kinetically trapped with methanol to deliver 3,4-disubstituted pyrrolidines 10 possessing three contiguous stereogenic centers. Experimentally, the irreversible formation of 8 was demonstrated since alkoxycyclization product 10 was not observed by its exposition to (PhO) 3 PAuBF 4 in the presence of methanol. This conclusion is in contrast with Fürstner's work, where related cyclobutenes 8 in the presence of a N-heterocyclic carbene 12 AuCl complex undergoes a rearrangement to give thermodynamically favored ring-expansion bicyclic [3.3.0] products through carbocation 11 [121]. Moreover, in contrast to the cationic stepwise mechanism commonly proposed [122,123], a recent work suggests a concerted pathway that precludes the intermediacy of vinyl-gold species 7 [124].
Catalysts 2020, 10, x FOR PEER REVIEW 5 of 49 demetalation, whereas carbocationic intermediate trans-9 is kinetically trapped with methanol to deliver 3,4-disubstituted pyrrolidines 10 possessing three contiguous stereogenic centers. Experimentally, the irreversible formation of 8 was demonstrated since alkoxycyclization product 10 was not observed by its exposition to (PhO)3PAuBF4 in the presence of methanol. This conclusion is in contrast with Fürstner's work, where related cyclobutenes 8′ in the presence of a N-heterocyclic carbene 12 AuCl complex undergoes a rearrangement to give thermodynamically favored ringexpansion bicyclic [3.3.0] products through carbocation 11 [121]. Moreover, in contrast to the cationic stepwise mechanism commonly proposed [122,123], a recent work suggests a concerted pathway that precludes the intermediacy of vinyl-gold species 7 [124]. [126]. Related phospha [6]helicenes have also shown high efficiency in terms of enantioselectivity and catalytic activity in this type of [2+2] cycloaddition reaction [127]. In general, switching from intramolecular to intermolecular processes implies drawbacks related to loss of regio-and stereoselectivity. However, similar conditions to those used in the intramolecular version allow for the maintaining of similar levels of regio-and enantioselectivity in gold-catalyzed intermolecular [2+2] cycloadditions of allenamides with alkenes [128][129][130][131]. In all the reported examples, the presence of an adjacent nitrogen in the allene moiety favors the regioselectivity by polarity induction generating, after coordination of the catalyst, a vinyl-gold conjugated acyliminium intermediate 20 that then evolves by nucleophilic addition of an electron-rich alkene (Scheme 5). In general, switching from intramolecular to intermolecular processes implies drawbacks related to loss of regio-and stereoselectivity. However, similar conditions to those used in the intramolecular version allow for the maintaining of similar levels of regio-and enantioselectivity in gold-catalyzed intermolecular [2+2] cycloadditions of allenamides with alkenes [128][129][130][131]. In all the reported examples, the presence of an adjacent nitrogen in the allene moiety favors the regioselectivity by polarity induction generating, after coordination of the catalyst, a vinyl-gold conjugated acyliminium intermediate 20 that then evolves by nucleophilic addition of an electron-rich alkene (Scheme 5).
Chen´s group first reported an efficient and selective intermolecular [2+2] cycloaddition approach to cyclobutane scaffolds using terminal allenamide 21 and alkenes 22 substituted by electron-donor groups, as enol ethers, in the presence of catalytic amounts of JohnphosAuCl/AgSBF 6 . Under these conditions, the corresponding products 23 with Z configuration are obtained (Scheme 6) [129].
In addition, the behaviour of ,-unsaturated N,N-alkyl hydrazones 26 in gold-catalyzed [2+2] cycloadditions with allenamides has been studied (Scheme 8). This transformation provides densely substituted cyclobutanes 27 with an all-carbon quaternary stereocenter. Excellent levels of regio-and diastereoselectivity and good yields are obtained, although in some cases mixtures of E/Z isomers were observed. Of note, the configuration of the alkene substrate is preserved in the final product, which is associated with steric factors [132]. In addition, the behaviour of α,β-unsaturated N,N-alkyl hydrazones 26 in gold-catalyzed [2+2] cycloadditions with allenamides has been studied (Scheme 8). This transformation provides densely substituted cyclobutanes 27 with an all-carbon quaternary stereocenter. Excellent levels of regio-and diastereoselectivity and good yields are obtained, although in some cases mixtures of E/Z isomers were observed. Of note, the configuration of the alkene substrate is preserved in the final product, which is associated with steric factors [132]. Gold complexes with phosphoramidite ligands 30-33 have shown satisfactory control of the enantioselectivity in reactions of N-allenylsulfonamides 28 with styrenes 29. The reaction works at low temperature due to the high reactivity of the allenes. The cycloaddition is compatible with electron withdrawing and electron donating substituents on the aromatic ring and no effect is observed with substitution at the different positions. In addition, the reaction conditions have allowed the synthesis of cyclobutanes containing challenging all-carbon quaternary stereocenters (Scheme 9) [133]. Gold complexes with phosphoramidite ligands 30-33 have shown satisfactory control of the enantioselectivity in reactions of N-allenylsulfonamides 28 with styrenes 29. The reaction works at low temperature due to the high reactivity of the allenes. The cycloaddition is compatible with electron withdrawing and electron donating substituents on the aromatic ring and no effect is observed with substitution at the different positions. In addition, the reaction conditions have allowed the synthesis of cyclobutanes containing challenging all-carbon quaternary stereocenters (Scheme 9) [133]. Gold complexes with phosphoramidite ligands 30-33 have shown satisfactory control of the enantioselectivity in reactions of N-allenylsulfonamides 28 with styrenes 29. The reaction works at low temperature due to the high reactivity of the allenes. The cycloaddition is compatible with electron withdrawing and electron donating substituents on the aromatic ring and no effect is observed with substitution at the different positions. In addition, the reaction conditions have allowed the synthesis of cyclobutanes containing challenging all-carbon quaternary stereocenters (Scheme 9) [133]. 3-(Propa-1,2-dien-1-yl)oxazolidin-2-one has been shown to be an efficient two carbon partner in a complete regio-and stereocontrolled intramolecular gold-catalyzed [2+2] cycloaddition with alkenes (Scheme 10). Trans isomer 35 is obtained regardless of alkene configuration. A stepwise cationic pathway involving cationic intermediates is proposed as mechanism. The nucleophilic attack of the alkene to a second cationic intermediate 34 would be the regioselectivity-determining step, favoring the formation of more stabilized benzylic or iminium cation [128,134]. Gold complex 36 with chiral 1,2,3-triazolylidene ligand has proven successful for enantioselective transformations over these substrates [135].
Catalysts 2020, 10, x FOR PEER REVIEW 9 of 49 3-(Propa-1,2-dien-1-yl)oxazolidin-2-one has been shown to be an efficient two carbon partner in a complete regio-and stereocontrolled intramolecular gold-catalyzed [2+2] cycloaddition with alkenes (Scheme 10). Trans isomer 35 is obtained regardless of alkene configuration. A stepwise cationic pathway involving cationic intermediates is proposed as mechanism. The nucleophilic attack of the alkene to a second cationic intermediate 34 would be the regioselectivity-determining step, favoring the formation of more stabilized benzylic or iminium cation [128,134]. Gold complex 36 with chiral 1,2,3-triazolylidene ligand has proven successful for enantioselective transformations over these substrates [135]. These intermolecular [2+2] cycloadditions of allenamides have been extended to the employment of indoles as olefin counterpart. In this regard, Bandini developed an enantioselective gold-catalyzed dearomative [2+2] cycloaddition of allenamides with indoles (Scheme 11a). This transformation enables direct access to methylenecyclobutane-fused indolines 37, featuring two consecutive quaternary stereogenic centers with excellent stereochemical control. The ring-closing event is favoured by the combined use of indoles carrying an electron withdrawing group at the N(1)position, which increases the electrophilicity of the dearomatized indolinine intermediate, and electron rich phosphines, which increase the nucleophilicity of alkenyl-gold species. DFT calculations support a polar non-concerted mechanism. Under kinetic conditions, tricyclic compound 38 is obtained, whereas isomeric cycloadduct 39 is generated under thermodynamic conditions (Scheme 11b). In both cases, the dearomatization process is the rate determining step [136,137]. These intermolecular [2+2] cycloadditions of allenamides have been extended to the employment of indoles as olefin counterpart. In this regard, Bandini developed an enantioselective gold-catalyzed dearomative [2+2] cycloaddition of allenamides with indoles (Scheme 11a). This transformation enables direct access to methylenecyclobutane-fused indolines 37, featuring two consecutive quaternary stereogenic centers with excellent stereochemical control. The ring-closing event is favoured by the combined use of indoles carrying an electron withdrawing group at the N(1)-position, which increases the electrophilicity of the dearomatized indolinine intermediate, and electron rich phosphines, which increase the nucleophilicity of alkenyl-gold species. DFT calculations support a polar non-concerted mechanism. Under kinetic conditions, tricyclic compound 38 is obtained, whereas isomeric cycloadduct 39 is generated under thermodynamic conditions (Scheme 11b). In both cases, the dearomatization process is the rate determining step [136,137].
Moreover, the reaction of 3-styrilindoles with N-allenamides grants enantioenriched cyclobutane derivatives and tetrahydrocarbazoles using H8-BINOL-derived phosphoramidite 40 and Me 2 SAuCl/AgNTf 2 . This transformation is highly dependent on the electronic nature of the indolic nitrogen. Enantioenriched non-fused cyclobutane derivatives are obtained with indoles bearing electron donating groups at the nitrogen, whereas tetrahydrocarbazoles are provided with electron withdrawing groups at that position. These results have been rationalized by DFT calculations (Scheme 12) [138]. Moreover, the reaction of 3-styrilindoles with N-allenamides grants enantioenriched cyclobutane derivatives and tetrahydrocarbazoles using H8-BINOL-derived phosphoramidite 40 and Me2SAuCl/AgNTf2. This transformation is highly dependent on the electronic nature of the indolic nitrogen. Enantioenriched non-fused cyclobutane derivatives are obtained with indoles bearing electron donating groups at the nitrogen, whereas tetrahydrocarbazoles are provided with electron withdrawing groups at that position. These results have been rationalized by DFT calculations (Scheme 12) [138]. Moreover, the reaction of 3-styrilindoles with N-allenamides grants enantioenriched cyclobutane derivatives and tetrahydrocarbazoles using H8-BINOL-derived phosphoramidite 40 and Me2SAuCl/AgNTf2. This transformation is highly dependent on the electronic nature of the indolic nitrogen. Enantioenriched non-fused cyclobutane derivatives are obtained with indoles bearing electron donating groups at the nitrogen, whereas tetrahydrocarbazoles are provided with electron withdr Xiang-Phos ligand 41 bearing two bulky adamantyl groups on the P atom gave the best results in the reaction of 3-styrylindoles with N-allenyl oxazolidinone (Scheme 13) [139]. Xiang-Phos ligand 41 bearing two bulky adamantyl groups on the P atom gave the best results in the reaction of 3-styrylindoles with N-allenyl oxazolidinone (Scheme 13) [139].
Alkene-Alkyne
The alkene-alkyne system has been extensively studied to generate cyclobutanes, cyclobutenes, and cyclobutanones through [2+2] cycloaddition in both intra-and intermolecular fashions [142]. 9-[158] enynes have been used as substrates affording a variety of different bicycles owing four-membered carbon rings. Most of these unsaturated systems required the use of bulky phosphines as ligands, mainly biarylphosphines. Some examples incorporating additionally 9-to 15-membered ring macrocycles has been also reported [159][160].
Regarding intramolecular [2+2] cycloaddition of enynes, alkene's addition to gold-activated alkynes afford cyclopropyl methyl carbenes as common intermediates in these processes. Several competing pathways arise, and their prevalence is regulated by the substitution pattern, catalysts, and conditions employed. The substrate may evolve through an exo-dig or endo-dig cyclization to render cyclobutenes. A general map for 1,6-enynes reactivity under gold catalysis is depicted in Scheme 17 [147,161,162]. A selective homodimerization of N-allenylsulfonamides to produce dialkylidencyclobutanes 50 occurs in good yields when alleneamide substrates are mixed with only 0.5 mol% of a gold catalyst. The reaction time decreases, and the yield increases with phosphite gold complexes as catalysts in combination with norbornene at 50 • C (Scheme 16) [130].
Allene-Allene
When a second allene molecule takes part in the cyclization process instead of an alkene, cyclobutanes with two exocyclic double bonds are achieved, which can be modified in further reactions.
Alkene-Alkyne
The alkene-alkyne system has been extensively studied to generate cyclobutanes, cyclobutenes, and cyclobutanones through [2+2] cycloaddition in both intra-and intermolecular fashions [142]. 9-[158] enynes have been used as substrates affording a variety of different bicycles owing four-membered carbon rings. Most of these unsaturated systems required the use of bulky phosphines as ligands, mainly biarylphosphines. Some examples incorporating additionally 9-to 15-membered ring macrocycles has been also reported [159][160].
Regarding intramolecular [2+2] cycloaddition of enynes, alkene's addition to gold-activated alkynes afford cyclopropyl methyl carbenes as common intermediates in these processes. Several competing pathways arise, and their prevalence is regulated by the substitution pattern, catalysts, and conditions employed. The substrate may evolve through an exo-dig or endo-dig cyclization to render cyclobutenes. A general map for 1,6-enynes reactivity under gold catalysis is depicted in Scheme 17 [147,161,162]. Scheme 16. Gold-catalyzed cyclodimerization of allenamides toward dialkylidencyclobutanes.
Alkene-Alkyne
The alkene-alkyne system has been extensively studied to generate cyclobutanes, cyclobutenes, and cyclobutanones through [2+2] cycloaddition in both intra-and intermolecular fashions [142]. 9-[158] enynes have been used as substrates affording a variety of different bicycles owing four-membered carbon rings. Most of these unsaturated systems required the use of bulky phosphines as ligands, mainly biarylphosphines. Some examples incorporating additionally 9-to 15-membered ring macrocycles has been also reported [159,160].
Regarding intramolecular [2+2] cycloaddition of enynes, alkene's addition to gold-activated alkynes afford cyclopropyl methyl carbenes as common intermediates in these processes. Several competing pathways arise, and their prevalence is regulated by the substitution pattern, catalysts, and conditions employed. The substrate may evolve through an exo-dig or endo-dig cyclization to render cyclobutenes. A general map for 1,6-enynes reactivity under gold catalysis is depicted in Scheme 17 [147,161,162]. The high energy activation from anti-52 to syn-52 (24.7 kcal/mol for R 1 = H and R 2 = Me) can be explained by the loss of conjugation between the gold carbene and the cyclopropane since there is a shortening of the cyclopropane and C=Au bonds, as well as a lengthening of the C-C bond connecting the cyclopropane and the gold carbene in the transition state [153]. Hence, such isomerization is rather unlikely under the reaction conditions. Nevertheless, an aryl group tethered to the carbene facilities rotation around the cyclopropane carbene bond by conjugation of the cyclopropane with the phenyl ring showing a low rotational barrier (8.6 kcal/mol) [147]. Direct formation of syn-52 can be achieved through an alternative pathway by syn-type attack of the alkene to the (alkyne)gold intermediate.
Syn-52 expands to form 54 with an activation energy of 11.9 kcal/mol. These results point to the anti to syn isomerization pathway possibly competing with the opening of anti-52. Afterwards, 54 can be obtained by a proton-elimination followed by protonolysis of the C-Au bond. 1,6-Enynes are reluctant to afford 53 due to the unstable configuration that represents an endo trans ring, thus isomerizing rapidly through 1,3-hydrogen migration to exo-54 analogue through acid catalysis [163,164].
Frequently, 1,3-dienes 55 are obtained as a result of a single-cleavage rearrangement, understood as a formal insertion of the alkyne between the carbons belonging to the alkene. Alternatively, a double-cleavage rearrangement is prone to happen, leading to products 56 in which both unsaturations have suffered C-C bond cleavage.
On the other hand, different cyclobutane derivatives can be obtained depending on the substrate substitution and the presence of other reagents on the reaction media. Thus, gold-catalyzed cycloisomerization of substituted 1,6-ene-ynamides 57 allows for access to cyclobutanones 60, through hydrolysis of the corresponding cyclobutenes initially formed, as reported by Cossy [145,148] and Yeh (Scheme 18) [152]. As trimethylsilyl-ynamides do not undergo fast protosilylation in the presence of AuCl, the mechanistic proposal involves the loss of the trimethylsilyl group by protonation of the double bond of the initially formed trimethylsilyl-substituted cyclobutene The high energy activation from anti-52 to syn-52 (24.7 kcal/mol for R 1 = H and R 2 = Me) can be explained by the loss of conjugation between the gold carbene and the cyclopropane since there is a shortening of the cyclopropane and C=Au bonds, as well as a lengthening of the C-C bond connecting the cyclopropane and the gold carbene in the transition state [153]. Hence, such isomerization is rather unlikely under the reaction conditions. Nevertheless, an aryl group tethered to the carbene facilities rotation around the cyclopropane carbene bond by conjugation of the cyclopropane with the phenyl ring showing a low rotational barrier (8.6 kcal/mol) [147]. Direct formation of syn-52 can be achieved through an alternative pathway by syn-type attack of the alkene to the (alkyne)gold intermediate.
Syn-52 expands to form 54 with an activation energy of 11.9 kcal/mol. These results point to the anti to syn isomerization pathway possibly competing with the opening of anti-52. Afterwards, 54 can be obtained by a proton-elimination followed by protonolysis of the C-Au bond. 1,6-Enynes are reluctant to afford 53 due to the unstable configuration that represents an endo trans ring, thus isomerizing rapidly through 1,3-hydrogen migration to exo-54 analogue through acid catalysis [163,164].
Frequently, 1,3-dienes 55 are obtained as a result of a single-cleavage rearrangement, understood as a formal insertion of the alkyne between the carbons belonging to the alkene. Alternatively, a double-cleavage rearrangement is prone to happen, leading to products 56 in which both unsaturations have suffered C-C bond cleavage.
Scheme 19. Gold-catalyzed tandem [2+2] cycloaddition/hydroarylation of enynes.
Propargylic esters 63, a particular type of 1,7-enynes, can be used as a mean to generate allenes 64 in situ, stemming from a syn 1,3-migration of the carboxylic ester group. Thus, alkylidencyclobutanes 65 are obtained through a formal [2+2] cycloaddition of allenenes, as reported by Zhang and Chang et al. [165][166][167]. Interestingly, the unsaturation thought to be activated by the gold catalyst is not the allene, but the alkene moiety to avoid unfavourable steric interactions (Scheme 20). Propargylic esters 63, a particular type of 1,7-enynes, can be used as a mean to generate allenes 64 in situ, stemming from a syn 1,3-migration of the carboxylic ester group. Thus, alkylidencyclobutanes 65 are obtained through a formal [2+2] cycloaddition of allenenes, as reported by Zhang and Chang et al. [165][166][167]. Interestingly, the unsaturation thought to be activated by the gold catalyst is not the allene, but the alkene moiety to avoid unfavourable steric interactions (Scheme 20).
Afterwards, a related enantioselective reaction that transforms indolyl substrates 66, bearing a propargyl ester unit in their structure, to complex tetracyclic scaffolds 67 has been carried out in the context of a study of structural features of acyclic diaminocarbene ligands by computational modelling (Scheme 21) [168]. A higher enantiodiscrimination is observed with bulkier alkyl substituents on alkyne group using gold catalyst 68. However, an opposite effect is detected with complex 69. In general, better yields and enantioselectivities are reached with complex 68 than with 69. Catalysts 2020, 10 Afterwards, a related enantioselective reaction that transforms indolyl substrates 66, bearing a propargyl ester unit in their structure, to complex tetracyclic scaffolds 67 has been carried out in the context of a study of structural features of acyclic diaminocarbene ligands by computational modelling (Scheme 21) [168]. A higher enantiodiscrimination is observed with bulkier alkyl substituents on alkyne group using gold catalyst 68. However, an opposite effect is detected with complex 69. In general, better yields and enantioselectivities are reached with complex 68 than with 69. On the other hand, the intermolecular version of [2+2] cycloaddition of alkynes with alkenes has been widely studied [169,170]. Potential problems with this approach lie in the competitive coordination of the alkene to the catalyst forming Au(I)-alkene complexes. These compounds lead to complex reaction mixtures or polymerizations in the presence of Au(I) complexes. Sterically hindered cationic Au(I) complexes minimize the competitive pathways.
Terminal alkynes are preferred due to the low reactivity associated to internal alkynes. Bulky Au(I)biphenylphosphine complexes form isolable σ,π-dicoordinated digold complexes 70 in the presence of phenylacetylene. These complexes catalyze an intermolecular [2+2] cycloaddition between phenylacetylene and α-methylstyrene with almost complete selectivity and higher yield for the cyclobutene compared to the corresponding mono Au(I) complex precursor (Scheme 22). In the presence of the latter, the Brønsted acid generated from the counteranion triggers α-methylstyrene dimerization and degradation of cyclobutene [170] On the other hand, the intermolecular version of [2+2] cycloaddition of alkynes with alkenes has been widely studied [169,170]. Potential problems with this approach lie in the competitive coordination of the alkene to the catalyst forming Au(I)-alkene complexes. These compounds lead to complex reaction mixtures or polymerizations in the presence of Au(I) complexes. Sterically hindered cationic Au(I) complexes minimize the competitive pathways.
On the other hand, the intermolecular version of [2+2] cycloaddition of alkynes with alkenes has been widely studied [169,170]. Potential problems with this approach lie in the competitive coordination of the alkene to the catalyst forming Au(I)-alkene complexes. These compounds lead to complex reaction mixtures or polymerizations in the presence of Au(I) complexes. Sterically hindered cationic Au(I) complexes minimize the competitive pathways.
Terminal alkynes are preferred due to the low reactivity associated to internal alkynes. Bulky Au(I)biphenylphosphine complexes form isolable σ,π-dicoordinated digold complexes 70 in the presence of phenylacetylene. These complexes catalyze an intermolecular [2+2] cycloaddition between phenylacetylene and -methylstyrene with almost complete selectivity and higher yield for the cyclobutene compared to the corresponding mono Au(I) complex precursor (Scheme 22). In the presence of the latter, the Brønsted acid generated from the counteranion triggers -methylstyrene dimerization and degradation of cyclobutene [170]. Later, it was shown that [tBuXPhosAu(MeCN)] + complex with the bulky and soft anion [BAr4 F ] -(3,5-bis(trifluoromethyl)phenylborate) improves the yields of [2+2] cycloaddition [171]. In this context, substituted cyclobutenes are synthesized in a regioselective fashion by intermolecular gold(I) catalyzed [2+2] cycloaddition with terminal electron-rich alkynes and aromatic or aliphatic alkenes. The mechanistic proposal based on kinetic studies and DFT calculations support as first intermediates cyclopropyl gold(I) carbenes and the involvement of an associative ligand exchange between the gold-coordinated alkene and the alkyne as the rate-limiting step [172]. The scope of the reaction includes reactions of 1,3-butadiynes 71 with alkenes 72. In these cases, a chemoselective [2+2] cycloaddition takes place only through the terminal alkyne to give alkynyl cyclobutenes 73, resulting from the coupling of the more substituted carbon of the alkene with the secondary carbon of the alkyne (Scheme 23) [172]. In this context, substituted cyclobutenes are synthesized in a regioselective fashion by intermolecular gold(I) catalyzed [2+2] cycloaddition with terminal electron-rich alkynes and aromatic or aliphatic alkenes. The mechanistic proposal based on kinetic studies and DFT calculations support as first intermediates cyclopropyl gold(I) carbenes and the involvement of an associative ligand exchange between the gold-coordinated alkene and the alkyne as the rate-limiting step [172]. The scope of the reaction includes reactions of 1,3-butadiynes 71 with alkenes 72. In these cases, a chemoselective [2+2] cycloaddition takes place only through the terminal alkyne to give alkynyl cyclobutenes 73, resulting from the coupling of the more substituted carbon of the alkene with the secondary carbon of the alkyne (Scheme 23) [172]. In this context, substituted cyclobutenes are synthesized in a regioselective fashion by intermolecular gold(I) catalyzed [2+2] cycloaddition with terminal electron-rich alkynes and aromatic or aliphatic alkenes. The mechanistic proposal based on kinetic studies and DFT calculations support as first intermediates cyclopropyl gold(I) carbenes and the involvement of an associative ligand exchange between the gold-coordinated alkene and the alkyne as the rate-limiting step [172]. The scope of the reaction includes reactions of 1,3-butadiynes 71 with alkenes 72. In these cases, a chemoselective [2+2] cycloaddition takes place only through the terminal alkyne to give alkynyl cyclobutenes 73, resulting from the coupling of the more substituted carbon of the alkene with the secondary carbon of the alkyne (Scheme 23) [172]. In the same way, 1-vinyl-3-substituted cyclobutenes 74 can be synthesized by reactions of alkenes with the corresponding 1,3-enynes [ Recently, [2+2] cycloaddition of chloroethynyls 75 with monosubstituted unactivated alkenes 76 has been described with excellent regioselectivities. The reaction is largely stereospecific with 1,2disubstituted unactivated alkenes (Scheme 25) [174]. The synthetic utility of the 1-chlorocyclobutene derivatives 77 has been demonstrated by their successful employment as substrates in cross-coupling reactions. Recently, [2+2] cycloaddition of chloroethynyls 75 with monosubstituted unactivated alkenes 76 has been described with excellent regioselectivities. The reaction is largely stereospecific with 1,2disubstituted unactivated alkenes (Scheme 25) [174]. The synthetic utility of the 1-chlorocyclobutene derivatives 77 has been demonstrated by their successful employment as substrates in cross-coupling reactions.
The first enantioselective example has been performed over di-and trisubstituted alkenes 78 with a gold catalyst bearing Josiphos ligand and [BAr4 F ] − counterion in order to reduce the amount of digold species formed (Scheme 26) [175]. Remarkably, this approach has been applied to the enantioselective total synthesis in 9 steps of Rumphellaone A, a terpenoid known for its cytotoxicity against human tumour cells. The first enantioselective example has been performed over di-and trisubstituted alkenes 78 with a gold catalyst bearing Josiphos ligand and [BAr 4 F ] − counterion in order to reduce the amount of digold species formed (Scheme 26) [175]. Remarkably, this approach has been applied to the enantioselective total synthesis in 9 steps of Rumphellaone A, a terpenoid known for its cytotoxicity against human tumour cells. The same reaction has been used to test a new atropisomeric teraryl monophosphine ligand, Joyaphos by Sparr et al. [176]. The results of this investigation have shown that (Sa)-Ph2JoyaphosAuCl and Cy2JoyaphosAuCl ( Figure 5) combined with AgSbF6 can be used to promote the desired reaction, although moderate yields and poor enantioselectivity are achieved. The same reaction has been used to test a new atropisomeric teraryl monophosphine ligand, Joyaphos by Sparr et al. [176]. The results of this investigation have shown that (Sa)-Ph 2 JoyaphosAuCl and Cy 2 JoyaphosAuCl ( Figure 5) combined with AgSbF 6 can be used to promote the desired reaction, although moderate yields and poor enantioselectivity are achieved. The same reaction has been used to test a new atropisomeric teraryl monophosphine ligand, Joyaphos by Sparr et al. [176]. The results of this investigation have shown that (Sa)-Ph2JoyaphosAuCl and Cy2JoyaphosAuCl ( Figure 5) combined with AgSbF6 can be used to promote the desired reaction, although moderate yields and poor enantioselectivity are achieved. Recently, Echavarren´s group has reported other enantioselective synthesis of Rumphellaone A, in 12 steps (ca. 8% yield), and Hushinone a norsesquiterpenoid found in the essential oils from the buds of Betula pubescens, in 16 steps (ca. 1.1 % yield) [160]. The key step is a diastereoselective gold(I)catalyzed [2+2] macrocyclization of 1,10-enyne 79 to build the cyclobutene moiety (Scheme 27). Recently, Echavarren's group has reported other enantioselective synthesis of Rumphellaone A, in 12 steps (ca. 8% yield), and Hushinone a norsesquiterpenoid found in the essential oils from the buds of Betula pubescens, in 16 steps (ca. 1.1% yield) [160]. The key step is a diastereoselective gold(I)-catalyzed [2+2] macrocyclization of 1,10-enyne 79 to build the cyclobutene moiety (Scheme 27). The same reaction has been used to test a new atropisomeric teraryl monophosphine ligand, Joyaphos by Sparr et al. [176]. The results of this investigation have shown that (Sa)-Ph2JoyaphosAuCl and Cy2JoyaphosAuCl ( Figure 5) combined with AgSbF6 can be used to promote the desired reaction, although moderate yields and poor enantioselectivity are achieved.
Alkyne-Alkyne
In contrast to 1,n-enynes, 1,n-diynes has been barely used as precursors of four-membered rings under gold catalysis. In fact, all the examples described entail the initial generation of an allene moiety that then evolve by formal [2+2] cycloaddition.
In this arena, alkynyl-propargylic pivaloates 80 have been shown to be suitable substrates to promote a gold-catalyzed [2+2] cyclization.
Alkyne-Alkyne
In contrast to 1,n-enynes, 1,n-diynes has been barely used as precursors of four-membered rings under gold catalysis. In fact, all the examples described entail the initial generation of an allene moiety that then evolve by formal [2+2] cycloaddition.
In this arena, alkynyl-propargylic pivaloates 80 have been shown to be suitable substrates to promote a gold-catalyzed [2+2] cyclization. Access to functionalized naphtho[b]cyclobutenes 85 with high stereoselectivity can be achieved by gold-catalyzed cascade cyclization of 1,7-diyn-3,6-bis(propargyl) carbonates 84. The cascade sequence involves a double 3,3-rearrangement forming bis(allenyl)carbonate 86. This is followed by a 6π-electrocyclic reaction to deliver a naphthyl derivative 87 that can be represented as a highly stabilized biradical 88 which by spontaneous cyclization affords cyclobutenyl dicarbonate 89. Finally, a decarbonylative cyclization provides 85 (pathway A). Alternatively, intramolecular nucleophilic attack of the allenic moiety on the gold-activated allene is proposed in pathway B affording oxocarbenium intermediate 90. Subsequent nucleophilic attack of the Au-C(sp 3 ) bond on the carbonyl moiety gives 89 which evolves in a similar sequence that in pathway A towards 85 (Scheme 30) [178].
On the other hand, intramolecular gold-catalyzed cycloisomerization of stable alkylidene tethered diynes 91 gives access to cyclobutene-fused azepines 92. A plausible mechanism has been proposed based on 1 H NMR studies (Scheme 31). Initially, nucleophilic addition of alkene to activated alkyne occurs through a 6-endo-dig attack supported by the nitrogen lone pair. The cleavage of C-N bond in 93 leads to an allylic cation 94 and consequent removal of the metal forms the allene 95. Finally, activation of the vicinal alkyne triggers a [2+2] cycloaddition with the former allene [179]. sequence involves a double 3,3-rearrangement forming bis(allenyl)carbonate 86. This is followed by a 6-electrocyclic reaction to deliver a naphthyl derivative 87 that can be represented as a highly stabilized biradical 88 which by spontaneous cyclization affords cyclobutenyl dicarbonate 89. Finally, a decarbonylative cyclization provides 85 (pathway A). Alternatively, intramolecular nucleophilic attack of the allenic moiety on the gold-activated allene is proposed in pathway B affording oxocarbenium intermediate 90. Subsequent nucleophilic attack of the Au-C(sp 3 ) bond on the carbonyl moiety gives 89 which evolves in a similar sequence that in pathway A towards 85 (Scheme 30) [178]. On the other hand, intramolecular gold-catalyzed cycloisomerization of stable alkylidene tethered diynes 91 gives access to cyclobutene-fused azepines 92. A plausible mechanism has been proposed based on 1 H NMR studies (Scheme 31). Initially, nucleophilic addition of alkene to activated alkyne occurs through a 6-endo-dig attack supported by the nitrogen lone pair. The cleavage of C-N bond in 93 leads to an allylic cation 94 and consequent removal of the metal forms the allene 95. Finally, activation of the vicinal alkyne triggers a [2+2] cycloaddition with the former allene [179]. [167]. A complete control of the product selectivity is gained by the study of the steric interactions between the alkyne moieties with the gold catalyst. Later, the same authors described the gold-catalyzed cycloisomerization of 1,6-diyne esters 98 to prepare chemoselectively bicyclo[3.2.0]hepta-1,5-dienes 99 (Scheme 32) [180]. Other examples of alkylidencyclobutene cores such as azabicyclo[4.2.0]octadienes 97 have been synthesized by Chan et al. from diyne substrates 96 through 1,3-migration/6-exo-dig cyclization/Prinstype [2+2]-cycloaddition [167]. A complete control of the product selectivity is gained by the study of the steric interactions between the alkyne moieties with the gold catalyst. Later, the same authors described the gold-catalyzed cycloisomerization of 1,6-diyne esters 98 to prepare chemoselectively bicyclo[3.2.0]hepta-1,5-dienes 99 (Scheme 32) [180]. Finally, an intermolecular approach has been reported by Shi et al. in which allenes are generated in situ from activated propargylic esters, using a silver-free gold catalyst.
[2+2] Cycloaddition works under mild conditions with high efficiency. The silver-free condition is crucial for the success of the transformation since the presence of a silver cation can activate the pyvaloate group as a leaving group, giving acyclic dimer 100 as the only observed product (Scheme 33). The [2+2] cycloaddition reaction is substrate-dependent and it was confirmed that it is a thermal reaction and, therefore, gold activation is not required for the cycloaddition event [181]. Finally, an intermolecular approach has been reported by Shi et al. in which allenes are generated in situ from activated propargylic esters, using a silver-free gold catalyst.
[2+2] Cycloaddition works under mild conditions with high efficiency. The silver-free condition is crucial for the success of the transformation since the presence of a silver cation can activate the pyvaloate group as a leaving group, giving acyclic dimer 100 as the only observed product (Scheme 33). The [2+2] cycloaddition reaction is substrate-dependent and it was confirmed that it is a thermal reaction and, therefore, gold activation is not required for the cycloaddition event [181].
Ring Expansion
A recurrent approach to access four-membered rings via gold(I) catalysis is the ring expansion of activated cyclopropane derivatives, taking advantage of the large ring strain associated to these carbocycles. This strategy to build the cyclobutane core typically implies the formation of a cyclopropylmethyl cation intermediate and a subsequent 1,2-alkyl migration. The noticeable πcharacter of the cyclopropyl group is responsible for its hyperconjugation, thus stabilizing these species as non-classical carbocations. In this context, vinylidene-, alkynyl-and alkylidencyclopropanes and vinyl-, alkynyl-and allenylcyclopropanols are suitable substrates to Scheme 33. Carbophilicity and oxophilicity competition in Au/Ag catalysis.
Ring Expansion
A recurrent approach to access four-membered rings via gold(I) catalysis is the ring expansion of activated cyclopropane derivatives, taking advantage of the large ring strain associated to these carbocycles. This strategy to build the cyclobutane core typically implies the formation of a cyclopropylmethyl cation intermediate and a subsequent 1,2-alkyl migration. The noticeable π-character of the cyclopropyl group is responsible for its hyperconjugation, thus stabilizing these species as non-classical carbocations. In this context, vinylidene-, alkynyl-and alkylidencyclopropanes and vinyl-, alkynyl-and allenylcyclopropanols are suitable substrates to react through a gold(I) catalyzed ring expansion to afford cyclobutane scaffolds. Considering this behavior, most of the examples reported can be categorized in pinacol-like or Wagner-Meerwin rearrangements.
Pinacol-Like Transformations
Cyclopropanols and related cyclopropyl ethers are suitable starting materials to produce cyclobutanes through gold-catalyzed ring expansion. Such a process relies on a 1,2-alkyl shift triggered by activation of a vicinal unsaturation. This alkyl migration is also assisted by the electron density donation from the lone pair of the oxygen.
In 2005, Toste's group reported the first examples of this kind of gold(I)-catalyzed ring expansion using 1-alkynylcyclopropanols 101 to render alkylidenecyclobutanones 103 in high yields and as single olefin isomers (Scheme 26) [182]. The process is tolerant of diverse substitution at the alkyne moiety and allows the employment of silyl ethers as substrates if 2 equivalents of methanol are added to the reaction media. In the mechanistic proposal, the coordination of cationic gold(I) to the alkyne induces a selective 1,2-alkyl shift of the most substituted chain in species 102, which is consistent with the experimental data. In this way, E isomer is obtained in a stereoselective and, remarkably, stereospecific manner regarding substituents on the ring. Related cycloalkylidenecyclobutanones 106 could be prepared from cyclopropanols 104 bearing a 1,6-diyne. This cascade process is proposed to entail a diyne cycloisomerization and the cyclopropyl expansion (Scheme 34) [183]. On the other hand, a single example of the preparation of an allyl cyclobutanone from a cyclopropanol bearing an allylic alcohol is included in a more general work that describes the ring expansion of 1,4-allylic diol derivatives under gold catalysis [184]. In the presence of a cationic gold complex and an oxidant such as a pyridine N-oxide, 1,3diketones can be constructed in a regioselective fashion from spirocyclic propargylic alcohols [185]. In the case of alkynylcyclopropanol 107, the corresponding -keto cyclobutanone 109 is obtained, although in moderate yield due to its instability upon column chromatography (Scheme 35). Notably, the oxidative cycloisomerization is completely selective and the formation or intermediacy of (E)-2benzylidenecyclobutanone is discarded. The plausible mechanism supported by NMR studies points out a favored coordination of gold to the N-oxide over coordination with the alkyne. The complex formed suffers oxidative addition at the alkyne moiety, leading to α-oxo gold-carbene species 108 which, after pinacol-like ring expansion, generates 1,3-diketone 109. In the presence of a cationic gold complex and an oxidant such as a pyridine N-oxide, 1,3-diketones can be constructed in a regioselective fashion from spirocyclic propargylic alcohols [185]. In the case of alkynylcyclopropanol 107, the corresponding β-keto cyclobutanone 109 is obtained, although in moderate yield due to its instability upon column chromatography (Scheme 35). Notably, the oxidative cycloisomerization is completely selective and the formation or intermediacy of (E)-2-benzylidenecyclobutanone is discarded. The plausible mechanism supported by NMR studies points out a favored coordination of gold to the N-oxide over coordination with the alkyne. The complex formed suffers oxidative addition at the alkyne moiety, leading to α-oxo gold-carbene species 108 which, after pinacol-like ring expansion, generates 1,3-diketone 109.
In the presence of a cationic gold complex and an oxidant such as a pyridine N-oxide, 1,3diketones can be constructed in a regioselective fashion from spirocyclic propargylic alcohols [185]. In the case of alkynylcyclopropanol 107, the corresponding -keto cyclobutanone 109 is obtained, although in moderate yield due to its instability upon column chromatography (Scheme 35). Notably, the oxidative cycloisomerization is completely selective and the formation or intermediacy of (E)-2benzylidenecyclobutanone is discarded. The plausible mechanism supported by NMR studies points out a favored coordination of gold to the N-oxide over coordination with the alkyne. The complex formed suffers oxidative addition at the alkyne moiety, leading to α-oxo gold-carbene species 108 which, after pinacol-like ring expansion, generates 1,3-diketone 109. More recently, O-substituted alkynylcyclopropyl allyl ethers 110 have also proven to be effective reactants for the construction of highly substituted alkylidencyclobutanones 113 using the same gold catalyst (Scheme 36) [186]. The reactions are accelerated by the presence of water and occur in moderate to good yields with broad scope. A thorough mechanistic study on the reported process has been accomplished including D and 18 More recently, O-substituted alkynylcyclopropyl allyl ethers 110 have also proven to be effective reactants for the construction of highly substituted alkylidencyclobutanones 113 using the same gold catalyst (Scheme 36) [186]. The reactions are accelerated by the presence of water and occur in moderate to good yields with broad scope. A thorough mechanistic study on the reported process has been accomplished including D and 18 Alkenylcyclopropanols and their ether derivatives are also precursors of cyclobutenones under gold catalysis. In an early work, Echavarren reported that alkenyl cyclopropyl ethyl ethers 114 in which the olefin is part of a terminal 1,6-enyne produce, in the presence of catalytic amounts of gold and water, variable mixtures of isomeric cyclobutane-fused tricyclic compounds 115 and 115a (Scheme 29) [187]. The isomeric ratio is mainly dependent on the catalyst employed: syn cycloadduct 115a is slightly favored in reactions conducted with cationic [JohnPhosAu(NCMe)]SbF6 catalyst, whereas anti tricyclic skeleton 115 is almost exclusively formed with AuCl.
This methodology has been employed as the key step of the total synthesis of Repraesentin F, using as starting material a cyclopropyl silyl ether 116 bearing an enyne system with an internal alkyne (Scheme 37). In this particular case, the higher selectivity to the desired tricyclic steroisomer is achieved with a cationic gold(I) catalyst derived from tBuXPhos biphenylphosphine ligand and with [BAr4 F ] -as counterion [188]. Alkenylcyclopropanols and their ether derivatives are also precursors of cyclobutenones under gold catalysis. In an early work, Echavarren reported that alkenyl cyclopropyl ethyl ethers 114 in which the olefin is part of a terminal 1,6-enyne produce, in the presence of catalytic amounts of gold and water, variable mixtures of isomeric cyclobutane-fused tricyclic compounds 115 and 115a (Scheme 29) [187]. The isomeric ratio is mainly dependent on the catalyst employed: syn cycloadduct 115a is slightly favored in reactions conducted with cationic [JohnPhosAu(NCMe)]SbF 6 catalyst, whereas anti tricyclic skeleton 115 is almost exclusively formed with AuCl.
This methodology has been employed as the key step of the total synthesis of Repraesentin F, using as starting material a cyclopropyl silyl ether 116 bearing an enyne system with an internal alkyne (Scheme 37). In this particular case, the higher selectivity to the desired tricyclic steroisomer is achieved with a cationic gold(I) catalyst derived from tBuXPhos biphenylphosphine ligand and with [BAr 4 F ] − as counterion [188].
115a is slightly favored in reactions conducted with cationic [JohnPhosAu(NCMe)]SbF6 catalyst, whereas anti tricyclic skeleton 115 is almost exclusively formed with AuCl. This methodology has been employed as the key step of the total synthesis of Repraesentin F, using as starting material a cyclopropyl silyl ether 116 bearing an enyne system with an internal alkyne (Scheme 37). In this particular case, the higher selectivity to the desired tricyclic steroisomer is achieved with a cationic gold(I) catalyst derived from tBuXPhos biphenylphosphine ligand and with [BAr4 F ] -as counterion [188]. However, studies accomplished by Voiturez's group on the enantioselective cyclization of related substrates 120 reveal a different behaviour. Under their optimized conditions that involve the use of a chiral bis(phosphine)digold(I) complex and wet toluene as solvent, selective evolution of the formation of cyclobutanones 121 occurs instead of the production of tricyclic adducts 115 (Scheme 39) [189]. The reactions take place with good yields and enantioselectivities although variable syn/anti diastereoselectivities are reached. However, studies accomplished by Voiturez's group on the enantioselective cyclization of related substrates 120 reveal a different behaviour. Under their optimized conditions that involve the use of a chiral bis(phosphine)digold(I) complex and wet toluene as solvent, selective evolution of the formation of cyclobutanones 121 occurs instead of the production of tricyclic adducts 115 (Scheme 39) [189]. The reactions take place with good yields and enantioselectivities although variable syn/anti diastereoselectivities are reached. A related reaction of substrate 122 having a cis-disubstituted cyclopropanol integrated in an 1,6enyne chain, gives bicyclic ketone 123 in excellent yield (Scheme 40) [190]. Cycloisomerization of cyclic olefin analogues 124 provides the corresponding tricyclic systems 125 in good yields and with total diastereoselection. Notably, the utility of this methodology for the rapid assembly of polycyclic ring systems is illustrated by its use in a key step of the total synthesis of the angular triquinane Ventricosene (Scheme 40). A related reaction of substrate 122 having a cis-disubstituted cyclopropanol integrated in an 1,6-enyne chain, gives bicyclic ketone 123 in excellent yield (Scheme 40) [190]. Cycloisomerization of cyclic olefin analogues 124 provides the corresponding tricyclic systems 125 in good yields and with total diastereoselection. Notably, the utility of this methodology for the rapid assembly of polycyclic ring systems is illustrated by its use in a key step of the total synthesis of the angular triquinane Ventricosene (Scheme 40). Efficient access to the tricyclic framework of protoilludanes 127 was described by Echavarren et al. through a related gold-catalyzed cycloisomerization of allene-alkenylcyclopropane derivatives 126 [191]. The process is also stereospecific, as demonstrates the different outcome of Z-and E-126 (Scheme 41). Efficient access to the tricyclic framework of protoilludanes 127 was described by Echavarren et al. through a related gold-catalyzed cycloisomerization of allene-alkenylcyclopropane derivatives 126 [191]. The process is also stereospecific, as demonstrates the different outcome of Zand E-126 (Scheme 41). Efficient access to the tricyclic framework of protoilludanes 127 was described by Echavarren et al. through a related gold-catalyzed cycloisomerization of allene-alkenylcyclopropane derivatives 126 [191]. The process is also stereospecific, as demonstrates the different outcome of Z-and E-126 (Scheme 41). On the other hand, analogous cyclopropyl enlargements to build the 4-membered ring are possible starting from 1-allenylcyclopropanols 128. In this sense, Toste reports the use of chiral dinuclear gold phosphine complexes for the construction of cyclobutanones 129 possessing a vinylsubstituted quaternary stereogenic center (Scheme 42) [192]. These reactions occur with good yields On the other hand, analogous cyclopropyl enlargements to build the 4-membered ring are possible starting from 1-allenylcyclopropanols 128. In this sense, Toste reports the use of chiral dinuclear gold phosphine complexes for the construction of cyclobutanones 129 possessing a vinyl-substituted quaternary stereogenic center (Scheme 42) [192]. These reactions occur with good yields and enantioselectivities, displaying broad scope and tolerance to functional groups. More recently, the same research group has combined visible light photoredox and gold catalysis to develop a novel approach to cyclic ketones from cycloalkanols through a ring expansion-Scheme 42. Enantioselective gold-catalyzed construction of cyclobutanones bearing a vinyl-substituted quaternary stereogenic center.
More recently, the same research group has combined visible light photoredox and gold catalysis to develop a novel approach to cyclic ketones from cycloalkanols through a ring expansion-oxidative arylation reaction in the presence of aryl diazonium salts. By this dual catalysis, using alkenyl or allenyl cycloalkanols 130, functionalized cyclobutanones 131 are furnished (Scheme 43) [193]. Mechanistic studies strongly suggest the initial formation of gold(III)-Ar species, by formal oxidative addition of aryldiazonium salts to gold(I), and the subsequent activation of the alkene or allene by this electrophilic gold(III) intermediate.
Wagner-Meerwein-Like Transformations
Polyunsaturated systems having a non-functionalized cyclopropane ring in their structure have also been proven as valuable precursors of cyclobutane containing cycloadducts in the presence of gold catalysts [75]. These processes typically involve an initial cycloisomerization upon metal activation of the alkyne and subsequent Wagner-Meerwein shift over the cyclopropylmethyl cation produced. In this sense, Min Shi's group described the gold(I)-catalyzed synthesis of cyclobutenefused carbazoles 133 by cycloisomerization of 1-(indol-3-yl)-3-alkyn-1-ols 132, a particular type of cyclopropyl embedded 1,5-enynes (Scheme 44) [194]. The authors proposed the key formation of a gold carbene intermediate 134 via nucleophilic attack of the indolyl group to the activated alkyne followed by 1,2-alkyl migration and posterior water elimination. Then, a ring expansion takes place that gives the tetracyclic skeleton and the final product after metal elimination.
Wagner-Meerwein-Like Transformations
Polyunsaturated systems having a non-functionalized cyclopropane ring in their structure have also been proven as valuable precursors of cyclobutane containing cycloadducts in the presence of gold catalysts [75]. These processes typically involve an initial cycloisomerization upon metal activation of the alkyne and subsequent Wagner-Meerwein shift over the cyclopropylmethyl cation produced. In this sense, Min Shi's group described the gold(I)-catalyzed synthesis of cyclobutene-fused carbazoles 133 by cycloisomerization of 1-(indol-3-yl)-3-alkyn-1-ols 132, a particular type of cyclopropyl embedded 1,5-enynes (Scheme 44) [194]. The authors proposed the key formation of a gold carbene intermediate 134 via nucleophilic attack of the indolyl group to the activated alkyne followed by 1,2-alkyl migration and posterior water elimination. Then, a ring expansion takes place that gives the tetracyclic skeleton and the final product after metal elimination.
produced. In this sense, Min Shi's group described the gold(I)-catalyzed synthesis of cyclobutenefused carbazoles 133 by cycloisomerization of 1-(indol-3-yl)-3-alkyn-1-ols 132, a particular type of cyclopropyl embedded 1,5-enynes (Scheme 44) [194]. The authors proposed the key formation of a gold carbene intermediate 134 via nucleophilic attack of the indolyl group to the activated alkyne followed by 1,2-alkyl migration and posterior water elimination. Then, a ring expansion takes place that gives the tetracyclic skeleton and the final product after metal elimination. The same group later reported that gold(I) catalysts are able to transform simple 1,5-enynes containing a cyclopropane ring in the alkyl chain into a variety of cyclobutane-fused cycloadducts depending on the substrate substitution, the temperature and the gold catalyst used [195,196]. Thus, reaction of (hetero)aryl-substituted 1,5-enynes 135 in dichloromethane at 0 °C using [JohnPhosAu(MeCN)]SbF6 as catalyst gives cyclobutane-fused 1,4-cyclohexadienes 136 in very high yields (Scheme 45). Moreover, the corresponding benzocyclobutenes 137 can be efficiently achieved by conducting the reactions under oxidative conditions. On the other hand, by simply controlling the reaction temperature and the gold(I) catalyst employed, three different products: biscyclopropanes 135, cyclobutane-fused 1,3-cyclohexadienes 139 and tricyclic cyclobutenes 140 can be selectively The same group later reported that gold(I) catalysts are able to transform simple 1,5-enynes containing a cyclopropane ring in the alkyl chain into a variety of cyclobutane-fused cycloadducts depending on the substrate substitution, the temperature and the gold catalyst used [195,196]. Thus, reaction of (hetero)aryl-substituted 1,5-enynes 135 in dichloromethane at 0 • C using [JohnPhosAu(MeCN)]SbF 6 as catalyst gives cyclobutane-fused 1,4-cyclohexadienes 136 in very high yields (Scheme 45). Moreover, the corresponding benzocyclobutenes 137 can be efficiently achieved by conducting the reactions under oxidative conditions. On the other hand, by simply controlling the reaction temperature and the gold(I) catalyst employed, three different products: biscyclopropanes 135, cyclobutane-fused 1,3-cyclohexadienes 139 and tricyclic cyclobutenes 140 can be selectively synthesized from 1,5-enynes 135, provided that an ortho-substituted arene is installed at their alkyne terminus (Scheme 45). In the same way, methylen-, vinyliden-, and alkylidencyclopropane enyne derivatives are also useful precursors for the construction of elaborated compounds containing a four-membered ring in their structure. This strategy was first illustrated in the transformation of 1,6-enyne 143 to fused cyclobutane 144 (Scheme 47) [190]. Interestingly, the introduction of aryl groups as alkyne substituents triggers a divergent evolution of the gold(I)-stabilized allyl cation intermediates to the diastereoselective construction of tetracyclic scaffolds 145 as a result of a final Nazarov-type electrocyclization (Scheme 47). In the same way, methylen-, vinyliden-, and alkylidencyclopropane enyne derivatives are also useful precursors for the construction of elaborated compounds containing a four-membered ring in their structure. This strategy was first illustrated in the transformation of 1,6-enyne 143 to fused cyclobutane 144 (Scheme 47) [190]. Interestingly, the introduction of aryl groups as alkyne substituents triggers a divergent evolution of the gold(I)-stabilized allyl cation intermediates to the diastereoselective construction of tetracyclic scaffolds 145 as a result of a final Nazarov-type electrocyclization (Scheme 47). The bicyclo[4.2.0]octane skeleton is also accessible from homologous methylidencyclopropyl 1,5-enynes 146 through a related mechanistic pathway. Thus, Gagné et al. described that these substrates undergo a 6-endo-dig cyclization followed by a Wagner-Merwein migration and 1,2hydrogen shift, over carbocation allylic species 147, to give bicyclo dienes 148 (Scheme 48) [197]. Good yields and moderate enantioselectivities are reached, being substitution at the cyclopropylidene moiety crucial for the enantioselection. Of note, Carreira et al. have further demonstrated the usefulness of this methodology in the total synthesis of a Harziane diterpenoid [198] (Scheme 48). Interestingly, also Gagné reported that particular 1,5-dienes containing a cyclohexane with an exocyclic cyclopropylidene are suitable substrates towards the construction of the bicyclo[4.2.0]oct- The bicyclo[4.2.0]octane skeleton is also accessible from homologous methylidencyclopropyl 1,5-enynes 146 through a related mechanistic pathway. Thus, Gagné et al. described that these substrates undergo a 6-endo-dig cyclization followed by a Wagner-Merwein migration and 1,2-hydrogen shift, over carbocation allylic species 147, to give bicyclo dienes 148 (Scheme 48) [197]. Good yields and moderate enantioselectivities are reached, being substitution at the cyclopropylidene moiety crucial for the enantioselection. Of note, Carreira et al. have further demonstrated the usefulness of this methodology in the total synthesis of a Harziane diterpenoid [198] (Scheme 48). Interestingly, also Gagné reported that particular 1,5-dienes containing a cyclohexane with an exocyclic cyclopropylidene are suitable substrates towards the construction of the bicyclo[4.2.0]oct-1-ene core [199]. In contrast, 1,5-enynes 146a bearing an additional alkynyl group directly attached to the methylencyclopropane unit, selectively produce tricyclic cycloadducts 151 as single diastereomers as a result of a complex cycloisomerization process of the initially formed bicyclic compounds 148 (Scheme 49) [201]. The corresponding 1,7-enynes possessing a methylenecyclopropane moiety are also suitable precursors of cyclobutene containing scaffolds through related mechanisms initiated by gold-alkyne activation. In this context, Min Shi's group has recently reported the preparation of cyclobutenyl substituted 1,2-dihydroquinolines 153 in moderate yields from aniline tethered 1,7-enynes 152 (Scheme 50) [202]. The process is limited to substrates bearing a terminal alkyne and aryl substituents at the methylidene unit. Moreover, reactions are not completely selective and the cyclobutenes are obtained accompanied with minor quantities of methylenecyclopropane isomers, which are [200]. The later reactions work nicely with aliphatic and electron deficient aromatic aldehydes, although it is limited to enynes with highly activated arenes at the alkyne.
In contrast, 1,5-enynes 146a bearing an additional alkynyl group directly attached to the methylencyclopropane unit, selectively produce tricyclic cycloadducts 151 as single diastereomers as a result of a complex cycloisomerization process of the initially formed bicyclic compounds 148 (Scheme 49) [201].
The corresponding 1,7-enynes possessing a methylenecyclopropane moiety are also suitable precursors of cyclobutene containing scaffolds through related mechanisms initiated by gold-alkyne activation. In this context, Min Shi's group has recently reported the preparation of cyclobutenyl substituted 1,2-dihydroquinolines 153 in moderate yields from aniline tethered 1,7-enynes 152 (Scheme 50) [202]. The process is limited to substrates bearing a terminal alkyne and aryl substituents at the methylidene unit. Moreover, reactions are not completely selective and the cyclobutenes are obtained accompanied with minor quantities of methylenecyclopropane isomers, which are selectively formed when using a silver catalyst. On the contrary, analogous ortho-(propargyloxy) aryl methylenecyclopropanes 154 (R = H) evolve through an intramolecular hydroarylation, instead of a 6-exo enyne cycloisomerization, followed by the ring enlargement to generate cyclobutenyl substituted 2H-chromenes 155 in good yields (Scheme 50) [203]. When blocking the reactive position of the arene to inhibit the hydroarylation process (R H) a new reaction occurs that efficiently provide cyclobutane fused dihydrobenzofuranes 156 (Scheme 50) [204]. This process is proposed to proceed via the initial ring expansion of the cyclopropane that generates a cyclobutene gold carbenoid species. This intermediate undergoes the nucleophilic attack of the oxygen followed by [2,3]-sigmatropic rearrangement and metal dissociation to give the observed tricyclic adducts. The process exhibits broad scope at both the arene and the alkyne moieties and only substrates bearing a terminal acetylene and bulky groups at the ortho position evolve through a different reaction pathway. Interestingly, the tricyclic adducts 156 could be also enantiomerically obtained using a chiral gold catalyst. In contrast, 1,5-enynes 146a bearing an additional alkynyl group directly attached to the methylencyclopropane unit, selectively produce tricyclic cycloadducts 151 as single diastereomers as a result of a complex cycloisomerization process of the initially formed bicyclic compounds 148 (Scheme 49) [201]. The corresponding 1,7-enynes possessing a methylenecyclopropane moiety are also suitable precursors of cyclobutene containing scaffolds through related mechanisms initiated by gold-alkyne activation. In this context, Min Shi's group has recently reported the preparation of cyclobutenyl substituted 1,2-dihydroquinolines 153 in moderate yields from aniline tethered 1,7-enynes 152 (Scheme 50) [202]. The process is limited to substrates bearing a terminal alkyne and aryl substituents at the methylidene unit. Moreover, reactions are not completely selective and the cyclobutenes are obtained accompanied with minor quantities of methylenecyclopropane isomers, which are selectively formed when using a silver catalyst. On the contrary, analogous ortho-(propargyloxy) aryl methylenecyclopropanes 154 (R = H) evolve through an intramolecular hydroarylation, instead of a 6-exo enyne cycloisomerization, followed by the ring enlargement to generate cyclobutenyl substituted 2H-chromenes 155 in good yields (Scheme 50) [203]. When blocking the reactive position of the arene to inhibit the hydroarylation process (R ≠ H) a new reaction occurs that efficiently provide cyclobutane fused dihydrobenzofuranes 156 (Scheme 50) [204]. This process is proposed to proceed via the initial ring expansion of the cyclopropane that generates a cyclobutene gold carbenoid species. This intermediate undergoes the nucleophilic attack of the oxygen followed by [2,3]-sigmatropic rearrangement and metal dissociation to give the observed tricyclic adducts. The process exhibits broad scope at both the arene and the alkyne moieties and only substrates bearing a terminal acetylene and bulky groups at the ortho position evolve through a different reaction pathway. Interestingly, the tricyclic adducts 156 could be also enantiomerically obtained using a chiral gold catalyst.
Scheme 50. Gold(I)-catalyzed ring expansion of aniline or benzyloxy tethered 1,7-enynes containing a cyclopropane unit towards cyclobutene scaffolds. Different heteropolycyclic scaffolds possessing a benzoazepine-fused cyclobutene core are furnished from related alkynylamide tethered methylenecyclopropanes 157 (Scheme 51) [205]. These reactions proceed via an initial 7-exo enyne cycloisomerization followed by the cyclopropylmethyl to cyclobutyl ring expansion that gives a tricyclic carbocationic intermediate 158, whose evolution is determined by the gold catalyst employed. Thus, deprotonation using PPh3AuCl/AgOTf triggers the Scheme 50. Gold(I)-catalyzed ring expansion of aniline or benzyloxy tethered 1,7-enynes containing a cyclopropane unit towards cyclobutene scaffolds. Different heteropolycyclic scaffolds possessing a benzoazepine-fused cyclobutene core are furnished from related alkynylamide tethered methylenecyclopropanes 157 (Scheme 51) [205]. These reactions proceed via an initial 7-exo enyne cycloisomerization followed by the cyclopropylmethyl to cyclobutyl ring expansion that gives a tricyclic carbocationic intermediate 158, whose evolution is determined by the gold catalyst employed. Thus, deprotonation using PPh 3 AuCl/AgOTf triggers the selective production of tricyclic compounds 159, whereas the use of JohnphosAuCl/NaBAr 4 F catalytic system promotes a selective intramolecular Friedel-Crafts type cyclization that finally delivers spirocyclic adducts 160. Interestingly, azepine-fused cyclobutanes with an aryl group at the bridgehead 162 could be prepared in an enantioselective manner from 1,6-enynes 161. These reactions occur with good yield and enantioselection although display limitations at the substitution at both the alkyne and olefin of the enyne (Scheme 52) [206]. The process involves asymmetric cyclopropanation of the methyliden unit, C-C cleavage and Wagner-Meerwein rearrangement. DFT calculations show that the chirality of the final product from the first cyclopropanation step is lost in the subsequent bond cleavage but is then regenerated in the Wagner−Meerwein rearrangement. Interestingly, azepine-fused cyclobutanes with an aryl group at the bridgehead 162 could be prepared in an enantioselective manner from 1,6-enynes 161. These reactions occur with good yield and enantioselection although display limitations at the substitution at both the alkyne and olefin of the enyne (Scheme 52) [206]. The process involves asymmetric cyclopropanation of the methyliden unit, C-C cleavage and Wagner-Meerwein rearrangement. DFT calculations show that the chirality of the final product from the first cyclopropanation step is lost in the subsequent bond cleavage but is then regenerated in the Wagner-Meerwein rearrangement. Interestingly, azepine-fused cyclobutanes with an aryl group at the bridgehead 162 could be prepared in an enantioselective manner from 1,6-enynes 161. These reactions occur with good yield and enantioselection although display limitations at the substitution at both the alkyne and olefin of the enyne (Scheme 52) [206]. The process involves asymmetric cyclopropanation of the methyliden unit, C-C cleavage and Wagner-Meerwein rearrangement. DFT calculations show that the chirality of the final product from the first cyclopropanation step is lost in the subsequent bond cleavage but is then regenerated in the Wagner−Meerwein rearrangement. Suitable functionalized vinylidencyclopropanes 163-165 are also valuable precursors of diverse cyclobutene-fused cycloadducts, in a similar fashion as the latter mentioned methylencyclopropanes. In this case, alkylidene cyclobutenyl carbene species 166, formed by ring expansion, are proposed as key intermediates and their evolution is governed by the substrate substitution (Scheme 53). Thus, vinylidencyclopropanes 163 having a pendant olefin, and a hydrogen or fluor substituent at C-4, selectively deliver benzocycloctane-fused cyclobutenes 167 via the intramolecular cyclopropanation of carbene intermediate 166a (Scheme 53a) [207]. Moreover, related aromatic substrates 164 furnish benzocycloheptane-fused cyclobutenes 168 through an intramolecular C-H insertion provided that a strong electron withdrawing group is present at the para position of the pendant arene (Scheme 53b) [208]. These processes display broad scope and, remarkably, their enantioselective versions have been also developed. However, competitive nucleophilic addition to the carbene intermediate are observed from most of the substrates substituted at C-4 and, therefore, no cyclobutane scaffolds are obtained with these substrates. Furthermore, intramolecular nucleophilic addition of oxygen nucleophiles to the carbene intermediate 166c yields methylene cyclobutanones 169 when vinylidencyclopropanes bearing both an alkyl substituent (R 1 ) and an aryl methyl ether are employed (Scheme 53c) [209]. Interestingly, both E and Z isomers 169 can be selectively obtained by controlling the gold ligand and the reaction conditions. In addition, alkylidene cyclobutanones are also accessible by reaction of unfunctionalized vinylidencyclopropanes with pyridine N-oxides via an intermolecular addition of the oxide to the alkylidene cyclobutenyl gold carbene species 166 initially generated [210]. (Scheme 53c) [209]. Interestingly, both E and Z isomers 169 can be selectively obtained by controlling the gold ligand and the reaction conditions. In addition, alkylidene cyclobutanones are also accessible by reaction of unfunctionalized vinylidencyclopropanes with pyridine N-oxides via an intermolecular addition of the oxide to the alkylidene cyclobutenyl gold carbene species 166 initially generated [210]. On the other hand, particular examples of gold-catalyzed synthesis of cyclobutene-fused compounds from cyclopropyl alkynes with a tethered oxygen nucleophile have been reported. In this sense, Zhang has described the preparation of densely functionalized bicyclo[3.2.0]heptanes 171 by a gold-catalyzed intermolecular reaction between alkynyl cyclopropyl ketones 170 and excess of ethyl vinyl ether and subsequent treatment with a Brønsted acid in the presence of water (Scheme 54) [211]. A mechanism has been proposed involving the initial cycloisomerization upon activation of the alkyne, followed by a 1,3-dipolar cycloaddition with the enol ether. Then, the construction of the cyclobutane is proposed to occur via a Wagner-Meerwein ring enlargement. Moreover, reactions of related 1-epoxi-1-alkynylcyclopropanes in the presence of water catalyzed by gold catalysts furnish functionalized bicyclic oxacyclic alcohols with good yields and diastereoselectivities [212]. Notably, the addition of halonium salts (NBS or NIS) allows the introduction of a halogen in the cyclization products [213]. Unactivated cyclopropyl alkynes can also be useful substrates for the preparation of cyclobutanes under gold catalysis provided that an appropriate nucleophile is present in the reaction media. Thus, reactions of a wide range of internal cyclopropyl alkynes 172 with sulfonamides in the presence of a cationic gold catalyst produce cyclobutanamines 174 in moderate to good yields [214]. A plausible mechanism for this transformation implies the alkyne activation by gold and subsequent ring expansion that generates a gold-stabilized allylic cyclobutyl carbocation 173, which is trapped by the sulfonamide to finally render the observed cyclobutyl sulfonamides after demetalation (Scheme 55). Interestingly, the same transformation is achieved using a magnetic nanoparticlesupported phosphine gold(I) complex [215]. External oxidants such as sulfoxides are also suitable reagents to trigger cyclobutane formation reactions from cyclopropyl alkynes via the well-documented formation of -oxo gold-carbene species. So, Liu developed an efficient gold(I)-catalyzed transformation of internal cyclopropyl alkynes 172a to cyclobutenyl ketones 176 using an excess of diphenyl sulfoxide as reactant (Scheme 56) [216]. This process begins with the selective intermolecular oxidation of the alkyne at the position respect to the cyclopropane thus producing -oxo gold-carbene intermediate 175, which then undergoes a Wagner-Merwein migration and posterior demetalation that accounts for the cyclobutene formation. An alternative concerted pathway in which the cyclopropyl expansion facilitates the cleavage of the O−S bond cannot be discarded. Unactivated cyclopropyl alkynes can also be useful substrates for the preparation of cyclobutanes under gold catalysis provided that an appropriate nucleophile is present in the reaction media. Thus, reactions of a wide range of internal cyclopropyl alkynes 172 with sulfonamides in the presence of a cationic gold catalyst produce cyclobutanamines 174 in moderate to good yields [214]. A plausible mechanism for this transformation implies the alkyne activation by gold and subsequent ring expansion that generates a gold-stabilized allylic cyclobutyl carbocation 173, which is trapped by the sulfonamide to finally render the observed cyclobutyl sulfonamides after demetalation (Scheme 55). Interestingly, the same transformation is achieved using a magnetic nanoparticle-supported phosphine gold(I) complex [215]. Unactivated cyclopropyl alkynes can also be useful substrates for the preparation of cyclobutanes under gold catalysis provided that an appropriate nucleophile is present in the reaction media. Thus, reactions of a wide range of internal cyclopropyl alkynes 172 with sulfonamides in the presence of a cationic gold catalyst produce cyclobutanamines 174 in moderate to good yields [214]. A plausible mechanism for this transformation implies the alkyne activation by gold and subsequent ring expansion that generates a gold-stabilized allylic cyclobutyl carbocation 173, which is trapped by the sulfonamide to finally render the observed cyclobutyl sulfonamides after demetalation (Scheme 55). Interestingly, the same transformation is achieved using a magnetic nanoparticlesupported phosphine gold(I) complex [215]. External oxidants such as sulfoxides are also suitable reagents to trigger cyclobutane formation reactions from cyclopropyl alkynes via the well-documented formation of -oxo gold-carbene species. So, Liu developed an efficient gold(I)-catalyzed transformation of internal cyclopropyl alkynes 172a to cyclobutenyl ketones 176 using an excess of diphenyl sulfoxide as reactant (Scheme 56) [216]. This process begins with the selective intermolecular oxidation of the alkyne at the position respect to the cyclopropane thus producing -oxo gold-carbene intermediate 175, which then undergoes a Wagner-Merwein migration and posterior demetalation that accounts for the cyclobutene formation. An alternative concerted pathway in which the cyclopropyl expansion facilitates the cleavage of the O−S bond cannot be discarded. External oxidants such as sulfoxides are also suitable reagents to trigger cyclobutane formation reactions from cyclopropyl alkynes via the well-documented formation of α-oxo gold-carbene species. So, Liu developed an efficient gold(I)-catalyzed transformation of internal cyclopropyl alkynes 172a to cyclobutenyl ketones 176 using an excess of diphenyl sulfoxide as reactant (Scheme 56) [216]. This process begins with the selective intermolecular oxidation of the alkyne at the β position respect to the cyclopropane thus producing α-oxo gold-carbene intermediate 175, which then undergoes a Wagner-Merwein migration and posterior demetalation that accounts for the cyclobutene formation.
An alternative concerted pathway in which the cyclopropyl expansion facilitates the cleavage of the O-S bond cannot be discarded.
External oxidants such as sulfoxides are also suitable reagents to trigger cyclobutane formation reactions from cyclopropyl alkynes via the well-documented formation of -oxo gold-carbene species. So, Liu developed an efficient gold(I)-catalyzed transformation of internal cyclopropyl alkynes 172a to cyclobutenyl ketones 176 using an excess of diphenyl sulfoxide as reactant (Scheme 56) [216]. This process begins with the selective intermolecular oxidation of the alkyne at the position respect to the cyclopropane thus producing -oxo gold-carbene intermediate 175, which then undergoes a Wagner-Merwein migration and posterior demetalation that accounts for the cyclobutene formation. An alternative concerted pathway in which the cyclopropyl expansion facilitates the cleavage of the O−S bond cannot be discarded. On the other hand, gold-catalyzed reaction of 4-methoxybut-2yn-1-ols 177 with two molecules of allyltrimethylsilane affords bicyclo[3.2.0]heptenes 181 in high diastereoselectity and good yields (Scheme 57) [217]. The proposed mechanism, based on deuterium-labeling experiments and the isolation of an intermediate, involves the formation of cyclopropyl alcohol 178 via propargylic substitution with an equivalent of allylsilane followed by enyne cycloisomerization. Then, the gold catalyst induces the formation of the cyclopropyl carbocation 179, which undergoes the ring enlargement to generate bicyclic carbocation 180 that gives the observed cycloadducts 181 by further reaction with the second unit of allylsilane. On the other hand, gold-catalyzed reaction of 4-methoxybut-2yn-1-ols 177 with two molecules of allyltrimethylsilane affords bicyclo[3.2.0]heptenes 181 in high diastereoselectity and good yields (Scheme 57) [217]. The proposed mechanism, based on deuterium-labeling experiments and the isolation of an intermediate, involves the formation of cyclopropyl alcohol 178 via propargylic substitution with an equivalent of allylsilane followed by enyne cycloisomerization. Then, the gold catalyst induces the formation of the cyclopropyl carbocation 179, which undergoes the ring enlargement to generate bicyclic carbocation 180 that gives the observed cycloadducts 181 by further reaction with the second unit of allylsilane.
Other Cyclization Approaches
Herein, other types of cyclizations from those shown previously are exposed. The first, reported by Li's group, involves the interaction of homopropargylic ethers 182 with pyridine N-oxide that gives -oxo gold-carbene intermediates 183. These species are intramolecularly trapped by the pendant alkoxy group to produce oxonium ylides species 184, whose rearrangement affords the observed -alkoxy cyclobutanones 185 in moderate yields (Scheme 58) [218]. The process is quite limited, as it only works with substrates bearing a terminal acetylene and electron-rich substituents at the homopropargylic position. Other substitution led to the formation of ,β-unsaturated carbonyl compounds [219], which are also obtained as side products with some of the suitable homopropargylic ethers 182.
Other Cyclization Approaches
Herein, other types of cyclizations from those shown previously are exposed. The first, reported by Li's group, involves the interaction of homopropargylic ethers 182 with pyridine N-oxide that gives α-oxo gold-carbene intermediates 183. These species are intramolecularly trapped by the pendant alkoxy group to produce oxonium ylides species 184, whose rearrangement affords the observed α-alkoxy cyclobutanones 185 in moderate yields (Scheme 58) [218]. The process is quite limited, as it only works with substrates bearing a terminal acetylene and electron-rich substituents at the homopropargylic position. Other substitution led to the formation of α,β-unsaturated carbonyl compounds [219], which are also obtained as side products with some of the suitable homopropargylic ethers 182.
gives -oxo gold-carbene intermediates 183. These species are intramolecularly trapped by the pendant alkoxy group to produce oxonium ylides species 184, whose rearrangement affords the observed -alkoxy cyclobutanones 185 in moderate yields (Scheme 58) [218]. The process is quite limited, as it only works with substrates bearing a terminal acetylene and electron-rich substituents at the homopropargylic position. Other substitution led to the formation of ,β-unsaturated carbonyl compounds [219], which are also obtained as side products with some of the suitable homopropargylic ethers 182. Scheme 58. Gold-catalyzed cyclization of homopropargylic ethers towards cyclobutanones.
Moreover, Hashmi's group employed the previously introduced σ,π-digold species for the construction of a couple of benzofused cyclobutane scaffolds 188 from thiophene embedded diynes 186 possessing a terminal alkyne at C-3 and a tert-butyl group at the internal alkyne in C-2 (Scheme Scheme 58. Gold-catalyzed cyclization of homopropargylic ethers towards cyclobutanones. Moreover, Hashmi's group employed the previously introduced σ,π-digold species for the construction of a couple of benzofused cyclobutane scaffolds 188 from thiophene embedded diynes 186 possessing a terminal alkyne at C-3 and a tert-butyl group at the internal alkyne in C-2 (Scheme 59) [220]. In this transformation, the gold catalyst coordinates the internal alkyne of the diyne system in a classical π-fashion whereas the terminal alkyne is σ-coordinated, thus producing the activated intermediate 187.
Then, the cyclobutane core is obtained through a 6-endo cyclization followed by a carbene C-H insertion.
Catalysts 2020, 10, x FOR PEER REVIEW 37 of 49 59) [220]. In this transformation, the gold catalyst coordinates the internal alkyne of the diyne system in a classical π-fashion whereas the terminal alkyne is σ-coordinated, thus producing the activated intermedia. Finally, the synthesis of cyclobutanones from alkynyl ketones through a gold-catalyzed oxidative process has been recently developed [221]. Thus, reaction of tert-butyl alkynyl ketone 189 with 8-isopropylquinoline N-oxide in the presence of a cationic gold catalyst gives -diketone--gold carbene intermediate 190, which then selectively evolves to observed cyclobutanones 191 by a C-H insertion (Scheme 60). Not surprisingly, variable selectivity to the desired cyclobutanones are achieved when more challenging substrates possessing different suitable hydrogens to undergo the C-H insertion are employed. For these substrates optimization of the gold catalyst is required. Finally, the synthesis of cyclobutanones from alkynyl ketones through a gold-catalyzed oxidative process has been recently developed [221]. Thus, reaction of tert-butyl alkynyl ketone 189 with 8-isopropylquinoline N-oxide in the presence of a cationic gold catalyst gives β-diketone-α-gold carbene intermediate 190, which then selectively evolves to observed cyclobutanones 191 by a C-H insertion (Scheme 60). Not surprisingly, variable selectivity to the desired cyclobutanones are achieved when more challenging substrates possessing different suitable hydrogens to undergo the C-H insertion are employed. For these substrates optimization of the gold catalyst is required. Catalysts 2020, 10, x FOR PEER REVIEW 38 of 49 Scheme 60. Synthesis of cyclobutanones from alkynyl ketones via oxidative gold catalysis.
Conclusions
The research recorded in this review highlights the great potential of gold catalysts for the construction of four carbon ring systems, a class of carbocycles widely present in natural products and very relevant intermediates in organic synthesis. Two main strategies have been developed in recent decades regarding on the reaction type: [2+2] cycloadditions and ring expansions. The progress of this field relies on the design and development of gold(I) catalysts to optimize yields, regio-, chemo-and stereoselectivities, as well as on the understanding of reaction mechanisms. These approaches allow to get complex cyclobutene derivatives under mild conditions, with high functional group tolerance and good yields, usually from easily available substrates. Enantioselective versions of some of these transformations employing gold complexes derived from chiral phosphine ligands have also been reported. In addition, examples of total synthesis of natural products following these methodologies have been shown to illustrate its relevance, advantages, and applications to access cyclobutane skeletons.
Conflicts of Interest:
The authors declare no conflict of interest. Scheme 60. Synthesis of cyclobutanones from alkynyl ketones via oxidative gold catalysis.
Conclusions
The research recorded in this review highlights the great potential of gold catalysts for the construction of four carbon ring systems, a class of carbocycles widely present in natural products and very relevant intermediates in organic synthesis. Two main strategies have been developed in recent decades regarding on the reaction type: [2+2] cycloadditions and ring expansions. The progress of this field relies on the design and development of gold(I) catalysts to optimize yields, regio-, chemoand stereoselectivities, as well as on the understanding of reaction mechanisms. These approaches allow to get complex cyclobutene derivatives under mild conditions, with high functional group tolerance and good yields, usually from easily available substrates. Enantioselective versions of some of these transformations employing gold complexes derived from chiral phosphine ligands have also been reported. In addition, examples of total synthesis of natural products following these methodologies have been shown to illustrate its relevance, advantages, and applications to access cyclobutane skeletons. | 17,105.2 | 2020-10-01T00:00:00.000 | [
"Chemistry"
] |
Study of Structure and Properties of Fe-Based Amorphous Ribbons after Pulsed Laser Interference Heating
The paper is devoted to the study of microstructural and magnetic properties of the Fe-based amorphous ribbons after interference pulsed laser heating. The ternary amorphous alloy FeSiB, as well as the multi-component alloys FeCuSiB and FeCuNbSiB, was subjected to laser pulses to induce crystallization in many microislands simultaneously. Structure and properties changes occurred in laser-heated dots. Detailed TEM analysis from a single dot shows the presence of FeSi(α) nanocrystals in the amorphous matrix. The FeSiB alloy is characterized after conventional crystallization by a dendritic structure; however, the alloys with copper as well copper and niobium additions are characterized by the formation of equiaxed crystals in the amorphous matrix. Amorphous alloys before and after the laser heating are soft magnetic; however, conventional crystallization leads to a deterioration of the soft magnetic properties of the material.
Introduction
Fe-based amorphous alloys have been extensively studied due to their very soft magnetic properties such as high saturation magnetization and near-zero coercivity (Ref 1-3). Due to their soft magnetic properties, these materials are used in production of magnetic cores, wires and shields. Since the amorphous structure of metallic glasses is metastable, crystallization processes of these alloys have been reported in many papers. It has been proved that the soft magnetic properties and the saturation magnetization of Fe-based amorphous materials can be improved by nanocrystallization of these materials (Ref 4,5). The best magnetic properties have been shown for alloys with nanometer size of grains ( Ref 1,4,6). That means, since a growing grain size leads to a deterioration of soft magnetic properties, the goal is to produce nanocrystalline alloys ( Ref 6). Unfortunately, the traditional heat treatment of FeSiB alloys leads to micrometer size of dendrites (Ref 3). On the other hand, annealing lead softens to a decrease in hardness and YoungÕs modulus, and consequently in a brittle crystalline state that makes the handling of the samples difficult ( Ref 7). A proper heat treatment that leads to nanocrystallization is difficult to achieve due to the simultaneous formation and growth of a crystalline phase which is characterized by harder magnetic properties than amorphous materials (Ref 8). Several unconventional techniques have been investigated to obtain nanostructured alloys ( . The pulsed laser interference process is based on the interference of at least two laser beams. The laser beams affect and create an interference pattern on the sample surface (Ref [15][16][17]. This process results in periodically arranged laser-heated microislands with changed structure and magnetic properties ( Ref 18,19).
The present study is focused on the microstructure evolution of Fe-based amorphous ribbons after pulsed laser interference heating and crystallization during annealing. The structure after laser and heat treatment is related to the magnetic properties of these alloys.
Experimental Materials and Methods
Fe-based amorphous alloys Fe 80 Si 11 B 9 , Fe 77 Cu 1 Si 13 B 9 and Fe 75 Cu 1 Nb 2 Si 1 3B 9 , with a thickness of 30 lm and a width of 25 mm and 10 mm, respectively, were studied. The chemical composition is well known as Metglass and Finemet alloys. Since these alloys are well investigated, it could well describe the influence of PLIH process on the structure and magnetic properties of these alloys. The amorphous alloys were fabricated by melt spinning (Ref 6,20). Pulsed laser interference heating (PLIH) was performed using a Nd:YAG (12 mm bar) laser with basic wavelength (1064 nm), 8-ns pulse time duration and 1 Hz frequency. The interference pattern was generated by a quartz tetrahedral prism with an apex angle of 172°. The device to perform PLIH process was described in (Ref 21). During the PLIH process, the number of consecutive laser pulses in the same area varied from 50 to 500, with 120 mJ energy. The laser beam parameters were selected to avoid ablation and cause maximum structure changes. Series of pulsed laser interference heating performed (Ref 19) and 120 mJ energy provided the most hopeful results. The samples were prepared by further heating the amorphous materials to 873 K at a rate of 20 K/min and slow cooling down to room temperature. Phase transformation of amorphous ribbon was tested using differential scanning calorimetry (DSC 800 Perkin Elmer) in a nitrogen atmosphere with 20 K/min heating rate. The structure of amorphous, laser-heated and annealed samples was studied using scanning (SEM-FEI Inspect S50) and transmission electron microscopy (TEM-Jeol JEM 2010-ARP). A Zeiss Libra 200 MC Cs STEM was used for structure characterization. The TEM samples were prepared by electropolishing of disks, cut from plain view. The magnetic properties were described by magnetic hysteresis loop measurements using a superconducting quantum interference device: (SQUID) magnetometer Quantum Design, MPMS, by applying an external field up to 4 T. The magnetic moment was determined with accuracy of 2% or better. The samples for SQUID measurements were cut as 3-mm disks and weighed to reference magnetization to the sample mass.
Phase Transformation
The DSC measurement of as-cast FeSiB amorphous ribbon indicates the presence of two exothermic peaks ( Fig. 1): First at 785 K (onset 775 K) was observed crystallization by the formation of primary BCC a-Fe(Si) phase, and second at 831 K (onset 825 K) it corresponds to the precipitation of intermetallic phases like Fe-B. It has been reported that at about 831 K, Fe 3 B intermetallic phase is formed ( Ref 2,22). It has been reported that copper addition decreases primary crystallization temperature (Ref 23). The alloys with copper and niobium addition were characterized by higher primary crystallization temperature ( Ref 24). The temperature of crystallization by annealing was selected 40 K above the peak of the precipitation intermetallic phase. Reason of this was crystallization after completed phase transformations.
Microstructure
Scanning electron microscopy (SEM) images showed a twodimensional structure composed of periodically arranged microislands fabricated by laser (Fig. 2). The distance between the microislands was about 16 lm. An increase in the number of laser pulses applied to the FeSiB alloy melted the surface and caused ablation of material ( Fig. 2 and 3). Irradiation with 500 consecutive laser pulses led to the formation of ripples in the laser-heated microislands ( Fig. 2d and 3b). According to Jia et al. (Ref 25), the creation of ripples is caused by the interference of the incident laser light and the scattered tangential light. The distance between the ripples is equal to the basic wavelength of Nd:YAG laser beam ($ 1064 nm) ( Ref 18).
Ripples were also observed at the FeCuSiB sample irradiated with 500 laser pulses (Fig. 3d) and at the FeCuNbSiB alloy after 200 laser pulses (Fig. 3e). The FeCuNbSiB alloy, laser heated with 500 laser pulses, was the alloy with the most heavily damaged surface in the microislands (Fig. 3f). This is probably due to thermophysical changes of alloys caused by Cu and Nb additions. Transmission electron microscopy (TEM) of the FeSiB laser-heated alloy showed nanocrystalline regions in the amorphous matrix (Fig. 4). Selected area electron diffraction (SAED) patterns can be attributed to the [À 103] zone axis of a-Fe(Si). The FeSiB amorphous alloy annealed at 873 K was characterized by a dendritic structure ( Fig. 5a and b), with a variable dendrite size from 100 to 500 nm. The SAED patterns indicate the occurrence of the a-Fe(Si) structure. Conventional crystallization of FeCuSiB and FeCuNbSiB alloys caused single nanocrystals in the amorphous matrix (Fig. 5c, d, e and f). Alloy additions (Cu and Nb) to FeSiB alloy resulted in crystal refining during annealing. The FeCuSiB alloy was characterized by crystals of a size from 50 to 100 nm, the FeCuSiB alloy by about 20 nm crystals size. As for the ternary alloy, the SAED patterns of these alloys indicate to presence of a-Fe(Si) structure.
High-Resolution Transmission Electron Microscopy
HRTEM of FeSiB alloys after laser irradiation with 120 mJ energy and 500 laser pulses showed a [111] a-Fe(Si) structure for one single nanocrystal (Fig. 6). At the boundary between the amorphous matrix and the crystals, a partial crystallization of the amorphous material can be observed (Fig. 7). The laser heating of the FeSiB alloy results in the formation of nanometer-sized crystals. The FeSiB sample after conventional crystallization by heating to 873 K shows a [001] a-Fe(Si) structure (Fig. 8). Annealing of this alloy resulted in large crystalline regions. Annealing of the alloy with copper addition led to the creation of a-Fe(Si) nanocrystals [111] oriented (Fig. 9a, b and c). The matrix was amorphous (Fig. 9d); however, partially crystallized material in the vicinity of crystals was observed (Fig. 9e). The alloy with copper and niobium additions contains the smallest size of nanocrystals. The HRTEM image shows nanocrystals with a size of a few nm (Fig. 10a). The structure of nanocrystals was a-Fe(Si) with [111] orientation in the HRTEM (Fig. 10b and c). Annealing of the FeSiB results in the creation of the highest size of crystals. The smallest size of crystals observed in alloys with copper and niobium addition can be explained according to Yoshizawa by the creation of Cu and Nb clusters which becomes nuclei for bcc Fe solid (Ref 1). The Fe grains were surrounded by Cu and Nb rich regions, which make grain growth difficult. Different treatments lead to the creation of different structures. The laser heating, as well as the annealing of the FeCuSiB and FeCuNbSiB alloys at 873 K, created nanocrystals in the amorphous matrix. The created structure was a-Fe(Si) [111] oriented. Conventional crystallization of the base alloy was with a dendritic structure composed of [001] a-Fe(Si) dendrites. The different orientation of dendrites may be caused by the preferred crystallographic orientation for dendritic growth which is [100] for a-Fe.
Magnetic Properties
To determine the magnetic properties of materials, magnetic hysteresis loop measurements were carried out. A FeSiB amorphous alloy shows the highest saturation magnetization at (Fig. 11). The FeCuSiB alloy has a 1.36 T saturation magnetization and the FeCuNbSiB alloy 1.3 T (Fig. 11b). The addition of copper and niobium leads to a decreased saturation magnetization of these alloys. Annealing at 873 K led to a decrease in the saturation magnetization of FeSiB and FeCuSiB alloys. The coercivity of amorphous alloys was 1.5 kA/m. After conventional crystallization, an increase in coercivity to 3 kA/m was observed. The FeCuNbSiB alloy after conventional crystallization was characterized by unchanged coercivity (Fig. 11c). The reason for this distinct behavior may be due to crystalsÕ size change; as Herzer (Ref 2) showed, increasing grain size led to increased coercivity of Fe-based nanocrystalline alloys, as observed here for FeSiB and FeCuSiB alloys. Niobium addition strongly decreases the size of crystals as discussed above: after crystallization, the crystals in FeCuNbSiB alloy were smaller than 20 nm (Fig. 10). The PLIH process led to a decreased saturation magnetization of these alloys except for the alloy with copper addition (Fig. 12). The FeSiB and the FeCuNbSiB alloys showed a saturation magnetization of 1.45 T and 1.27 T, respectively, after PLIH with 120 mJ and 300 laser pulses. The FeCuSiB alloy after laser heating was characterized by a saturation magnetization that was 0.13 T higher than for the amorphous alloy (Fig. 12b). The coercivity of 1.5 kA/m was unchanged for these materials after the laser heating process (Fig. 12c).
Conclusions
The application of the pulsed laser interference heating process formed two-dimensional structures. SEM images showed periodically arranged laser-heated microisland. The crystalline a-Fe(Si) structures were observed in these microislands. The matrix remained amorphous. HRTEM of the FeSiB alloy after the PLIH process showed single a-Fe(Si) nanocrystals and a partially crystallized amorphous matrix. Similar structures were obtained by Wu et al. (Ref 12) as well as Katakam et al. (Ref 13). It allows concluding that by the focused beams it is possible to obtain nanostructures, but only in laser-heated microisland. The crystallization during anneal- The SQUID measurement showed very soft magnetic properties of the amorphous ribbons. The materials that were crystallized during annealing were characterized by a deterio-rated saturation magnetization and coercivity, except for the FeCuNbSiB alloy. This could be correlated with the structure of these alloys. FeSiB alloy and FeCuSiB alloy were characterized by coarse dendrite/crystals structure, but FeCuNbSiB alloy had a nanocrystalline structure. The laser heating does not influence to the coercivity of these alloys, but it led to lower saturation magnetization for the FeSiB and FeCuNbSiB alloys. The volume of laser crystallized material must be not enough to increase the soft magnetic properties of these alloys. The FeCuSiB alloy showed higher saturation magnetization after the PLIH process than before the laser heating. Katakam et al. The laser beam energy is delivered at only 8 ns and is not enough to heat the higher volume of material, so the structural changes are observed only in laser-affected areas. Magnetic properties correlate with a structure so insufficient evolution of the structure will not change the magnetic properties. Perhaps the higher supplied energy to the material will increase the soft magnetic properties. It could be realized by increasing the laser beam energy or the number of laser pulses, as well as the time of laser pulse. The higher number of laser pulses, as well as the higher laser beam energy [our previous investigation (Ref 19)], leads to undesirable ablation. Changes of the laser pulses time could provide different structures, but it should be proved.
Open Access
This article is licensed under a Creative Commons Attribution 4.0 International License, which permits use, sharing, adaptation, distribution and reproduction in any medium or format, as long as you give appropriate credit to the original author(s) and the source, provide a link to the Creative Commons licence, and indicate if changes were made. The images or other third party material in this article are included in the article's Creative Commons licence, unless indicated otherwise in a credit line to the material. If material is not included in the article's Creative Commons licence and your intended use is not permitted by statutory regulation or exceeds the permitted use, you will need to obtain permission directly from the copyright holder. To view a copy of this licence, visit http://creativecommons.org/licenses/by/4.0/. | 3,278.8 | 2020-09-15T00:00:00.000 | [
"Materials Science"
] |
Multilateral Nuclear Approach to Nuclear Fuel Cycles
The Fukushima Daiichi Nuclear Power Plant (Fu-NPP) accident, which occurred in March 2011, has significantly influenced the recent trend in the growing interest in nuclear reactor deployment. On the other hand, it is common knowledge that there remain problems in global climate changes and energy security in the long term view that human-beings are ob‐ liged to solve. Meanwhile, the issue on treatment of spent fuels (SFs) has become remarka‐ bly recognized. These facts should result in the continuous needs of nuclear energy and proper management of SFs, i.e., well-organized nuclear fuel cycle (NFC) services, including uranium mining, refining, conversion, enrichment, reconversion, fuel fabrication, spent fuel treatment such as storage, reprocessing, and repository.
Introduction
The Fukushima Daiichi Nuclear Power Plant (Fu-NPP) accident, which occurred in March 2011, has significantly influenced the recent trend in the growing interest in nuclear reactor deployment.On the other hand, it is common knowledge that there remain problems in global climate changes and energy security in the long term view that human-beings are obliged to solve.Meanwhile, the issue on treatment of spent fuels (SFs) has become remarkably recognized.These facts should result in the continuous needs of nuclear energy and proper management of SFs, i.e., well-organized nuclear fuel cycle (NFC) services, including uranium mining, refining, conversion, enrichment, reconversion, fuel fabrication, spent fuel treatment such as storage, reprocessing, and repository.
The concerns for the nuclear proliferation of so-called "Sensitive Nuclear Technologies (SNTs)", and weapon-use nuclear materials, namely, enrichment technology (frontend), spent fuel reprocessing technology (backend) and fissile materials will increase.The latter includes concern with worldwide increase in the amount of SFs, which may have to be stored in individual states.Namely, there will be growing concerns from the nuclear non-proliferation and security perspectives that plutonium may globally proliferate as a form of SFs.
Measures for nuclear non-proliferation have so far been taken mainly by the combination of institutional systems and supply side approaches (see Fig. 1).International society has been responding to the above-mentioned concerns by strengthening schematic measures centered around the safeguards under the NPT and Convention on the Physical Protection of Nuclear Material, etc. Bilateral agreements represent the latter one (supply side approach) particularly those between the US and individual states that have been functioning strongly.However, increase in the supply of fuel source materials from the Eastern Block has been remarkable in recent years, as shown in Fig. 2, which may potentially weaken the influence of the Western Block on nuclear non-proliferation.
The measures for enhancement of nuclear non-proliferation on the supply side that mainly consist of the nuclear power technology advanced countries may interfere with the inalienable right of peaceful uses of nuclear power that is guaranteed by Article 4 of the NPT.Thus, there is a need to develop nuclear non-proliferation measures with high non-proliferation capacity based on new concepts which are completely different from the conventional ones.In addition, as for the nuclear security for handling SNTs and nuclear materials as well as safety management of nuclear facility operations, the conventional state-by-state efforts have limitations from the viewpoints of effectiveness, efficiency, and economic reasonability.Demand side approach represented by Multilateral Nuclear Approach (MNA), where services on the frontend and the backend are provided to the states possessing nuclear power plants without interfering with the inalienable right in NPT and measures for nuclear nonproliferation properly function, may be one of the most effective and efficient goals to solve all the problems discussed above.Originally the MNA was proposed as an idea to reduce the possibility of nuclear proliferation of sensitive technologies by supplying enrichment and reprocessing services to newcomer countries [1].
Regional MNA, e.g., in Asian regions, may also complement or reinforce the weakened nonproliferation regime of Western Block-based regions.In the foreseeable future, most of newcomer countries utilizing nuclear energy would like to have a reliable system of fresh fuel supply and spent fuel treatment services, free of any political disruptions to fuel their nuclear reactors.
Several proposals [2,3] on the Multilateral Nuclear Approach (MNA) have recently been studied and a few are now ready to be implemented, in which no restraint of the peaceful use of nuclear energy due to the issues of proliferation of sensitive technologies is taken into account.Recent discussions, however, tend to focus on reliable fuel supply, namely frontend of NFC, where proliferation of uranium enrichment can be deterred.At the same time, the MNA capability to provide assurance/service that the SFs be managed properly is actually more important [4].
Storing SFs as well as possessing those in power reactors in individual countries remains problematic not only for Safety but also the risks in nuclear proliferation/Safeguards and Security (3S), due to the presence of large amounts of imbedded plutonium (Pu).Although Pu in SFs is protected by its high radiation dose rate, the technology to separate the plutonium from SFs (reprocessing) is not as difficult as uranium enrichment technology.It is therefore important to address the issues associated with the backend of the nuclear fuel cycle and to propose to properly manage/treat SFs.In this context, MNA may also be beneficial from the viewpoints of 3S.
Historical review of international framework [5]
From the perspective of preventing proliferation of SNTs, the concept of "international control" has been proposed for a long time.
The old one is the international control of nuclear materials, which was developed under the Truman Administration in 1946 (i.e.pooling all nuclear materials, etc. in an international organization and lending them to states that want them).In January 1946, the United Nations Atomic Energy Commission (UNAEC) was founded based on the proposal by the US, the UK and Canada in November 1945 [6] that asked the international control of atomic energy, i.e., "control of atomic energy to the extent necessary to ensure its use only for peaceful purposes" and "elimination from national armaments of atomic weapons and of all other major weapons adaptable to mass destruction".In this way, MNA has been encouraged in line with the use of nuclear energy.
In the US, "A Report on the International Control of Atomic Energy" (also known as the "Acheson-Lilienthal Report") was prepared for discussion in the UNAEC.Based on the Declaration on the Atomic Bomb, the Report proposed to establish a new international organization called the "Atomic Development Authority (ADA)" which owns all fissionable material and controls them under effective leasing arrangements [7].The report also proposed that the ADA would be in charge of all "dangerous" activities relating to raw materials, construction and operation of production plants, and the conducting of research into explosives, while "non-dangerous" nuclear activities, such as the construction and operation of power-producing piles, "would be left in national hands".It is interesting to note that in the report, all nuclear fuel cycle activities, except nuclear reactors, were categorized as "dangerous activities" and should not be conducted by an individual state.
In 1946, Bernard Baruch, the US representative to the UNAEC, submitted his plan for the international control of nuclear energy based on the Acheson-Lilienthal Report [8].However, he modified the report by inserting the prohibition of the development of nuclear-weapons capability by new states and punishment for violations of such prohibition.The plan was not accepted by the Soviet Union and as a result, the international control of nuclear energy did not bear fruit in the 1940s.
This plan was later put on the table of the UNAEC in the form of "Baruch Plan" by UN Representative B. Baruch.However, the Plan did not take off successfully because it was contradicting with the US's free enterprise system of that time as it was promoting international ownership of the US technology.It also reached a deadlock in the negotiations between the US and the Soviet Union.However, the Plan triggered the "Age of International Collaboration for Peaceful Use of Nuclear Energy" in the "Atoms for Peace" speech by US President Eisenhower in 1953 at the UN.In this initiative, the uranium bank (reserve) with an intension of international management of fissile materials was proposed.After these debates, the International Atomic Energy Agency (IAEA) was established in 1957.Provision of nuclear materials, etc. became one of the missions of the IAEA.However, the uranium bank plan was eventually abandoned because a) uranium supply was not as limited as was initially envisioned, and b) competition of commercial nuclear energy technology/supply of nuclear materials in the major supplying states based on the above speech was intensified.
In post-war Europe, the European Atomic Energy Community (EURATOM) was established to promote nuclear energy development.The most important requirement of the EURA-TOM Convention was "to guarantee nuclear materials supply" by the member states [2].At the same time, the Convention had safeguard systems to ensure that the nuclear materials within EURATOM were to be used only for peaceful uses.
International debate with regards to exporting nuclear technology and material/equipment promoted another international framework concerning the supply.In 1971, the Zangger Committee was established.The member states shall apply the IAEA's safeguards to the exported "nuclear materials" when exporting them to the non-NPT member states without nuclear weapons as well as when transporting them from these nonnuclear weapons states.The Committee also created a list of equipment as subjects of the regulation.Meanwhile, after the first nuclear test by India, the Nuclear Suppliers Group (NSG) was established in 1974 for a similar purpose.The NSG controls exports based on the so-called "NSG Guidelines", the guidelines designed for the states which export nuclear energy related equipment, material and technologies (it is a "gentleman's agreement" without any legal binding power).
In 1975, the IAEA began the exploration of the first Regional Nuclear Fuel Cycle Center (RFCC) [9] and assessed the advantages of applying backend to the RFCC.The RFCC report examined and presented basic research from international and regional approaches regarding the backend of fuel cycle in various geographical sites.From 1977 to 1980, the International Nuclear Fuel Cycle Evaluation (INFCE) [10] was conducted, and the effectiveness of nuclear fuel cycle was thoroughly evaluated by 8 working groups (WGs).Through this activity, many WGs picked up "fuel cycle center" and described it as a systematic arrangement to strengthen nuclear non-proliferation.Furthermore, for the spent fuel issues, they considered the fuel cycle as a solution that includes legal framework and multinational arrangement.Based on the results of the INFCE, the IAEA supported the experts group to examine the concept of international plutonium storage (IPS) [9], established the Committee for Assurance of Supply (CAS) [9] in 1980 and continued the deliberation until 1987.The experts' examination concluded that the multilateral approach was technically and economically feasible but there were still issues in terms of difficulty in prerequisites for participation and transfer of rights towards nuclear non-proliferation.Most of those activities, initiated by the US-initiatives, could not agree on the non-proliferation commitments and conditions that would entitle states to participate in the multilateral activities" [11], because nuclear developed states in Western Europe and Japan had already engaged in the development of their own sensitive capabilities.Therefore, they tried to maintain their activities, and not let them be interfered with by such initiatives.Developing states, especially NAM states, argued that any requirements for the non-proliferation commitment of not engaging in sensitive nuclear activities were against Article IV of the NPT.They also insisted such a requirement would discriminate between the "haves" and "have-nots" of sensitive capabilities, in addition to there being an existing discrimination by the NPT between "NWS" and "NNWS".The US was thus left alone and could not gain enough support for promoting its initiatives any further.Together with Cold War tensions, a decline in the growth rate of the US economy and a decrease of energy demand due to the second oil shock in 1979, and discouragement of nuclear energy use after TMI and the Chernobyl accidents in 1979 and 1986, the US itself lost its motivation for MNA.
At GLOBAL 93, an international conference, the "International Monitored Retrievable Storage System (IMRSS)" [12] was proposed by Dr. Häfele from Germany.IMRSS proposed that spent fuel and plutonium shall be stored in a retrievable condition under monitoring by an international entity.It chose the IAEA as a desirable entity to lead the initiative.Although it was considered a temporally measure to buy some time until the conclusion of whether SFs would be directly disposed or plutonium would be retrieved, there was no development thereafter.Dr. A. Suzuki of the University of Tokyo made a proposal for spent fuel storage in the East Asia region, and Dr. J-S.Choi of CISAC/Stanford University made a proposal for the regional treaty including regional spent fuel storage.Their proposals show the significance of the systems in which the host states offer interim storage of SFs for a limited time (40 to 50 years), even though the handling of SFs from other states is not easy.
In 1994, the US and Russia agreed that the US would purchase 500 tons of highly-enriched uranium (HEU) from Russia, convert it to low-enriched uranium (LEU) and make peaceful uses of it.Furthermore, both states agreed that each state would declare 50 tons of excess plutonium to be used for defense purposes, dismantle and retrieve 34 tons from nuclear weapons, and convert it to power generating fuel as MOX.For the purpose of nuclear non-proliferation, the US also began the "Foreign Research Reactor Spent Nuclear Fuel Acceptance Program (FRRSNFA)" in 1996 to accept the US-origin spent HEU and LEU fuels from foreign research reactors by May 2009.Furthermore, under the Russian Research Reactor Fuel Return (RRRFR) Program, some 2 tons of HEU and some 2.5 tons of LEU SFs, which were previously supplied by Soviet Union/Russia to foreign reactors, were shipped to the Mayak reprocessing complex near Chelyabinsk.The US and the Russian Federation cooperated in several repatriation projects for Russian-origin HEU fuels.
Based on the recognition that SFs and high level waste (HLW) are the common critical issues which could be factors to hinder nuclear energy promotion in the East Asia region, the Pacific Nuclear Council (PNC) began deliberation to promote understanding and cooperation for the management of SFs and HLW among the PNC members and to investigate possibilities of the International Interim Storage Scheme (IISS) in 1997.The IISS is managed at national, regional, or international levels and is to augment (not to replace) the national system.The IISS operates during the contract period from the time when SFs and HLW are deposited to the storage facility in the host state till the time when "they are returned to the originating state".The host state would be responsible for safety and safeguards of the storage facility and receive financial compensation from the contact member state, which is the owner of the SFs and HLW.
In reality, the interim storage of SFs, a part of a reprocessing contract, had been offered by reprocessing operators such as the BNFL and the AREVA.With this system, the state which makes a reprocessing contract can store SFs as long as it is stored in the reprocessing facility; however the separated plutonium and HLW at the time of reprocessing would be returned to the state.On the other hand, the concepts of the IMRSS, the RSSFEA, regional treaty and the IISS demand the host state to store or dispose of other state's SF.However, this is not easy in reality.
Recent proposalsp [13, 14]
The concerns about nuclear proliferation by states and the acquisition of nuclear weapons by terrorists has grown after nuclear testing by India/Pakistan in 1998 and the terrorist attack on September 11, 2001.The nuclear weapons black market network issues by Democratic People's Republic of Korea (DPRK, hereafter referred to as North Korea), Libya, Iran and A.Q. Khan are driving the international society to make efforts through various trials and proposals in preventing proliferation of the SNT related to fuel cycle such as isotope separation and reprocessing.
The proposals made by Ex-Director General of the IAEA, Dr. M. ElBaradei, in October 2003 presented that (1) reprocessing and enrichment operations must be restricted under the multinational control, (2) nuclear energy system shall have nuclear non-proliferation resistance, and (3) multinational approaches shall be considered for the management and disposal of SFs and radioactive waste.However, it was anticipated that his idea of a multilateral system of SNTs and radioactive substances would take a long time to overcome the issues.
Former US President G.W. Bush strongly demanded in his speech at the Defense University in February 2004 that exporting SNTs should be limited to the states which were already using them on a full scale and respecting the Additional Protocol.However, this proposal may lead to international cartels and may split the member states into the states with SNTs and without SNTs.The "Nuclear Fuel Leasing" proposal by V. Rice, et al. and "Nuclear Fuel Service Assurance Initiative" proposal by E. Moniz, et al. expect the improvement of nuclear non-proliferation though institutionalization.However, the proposals still contain a concern over supply assurance to the user states as well as a concern over the dichotomization of the member states, similar to the other proposals.
Later, a group of experts for multinational nuclear (fuel cycle) approaches (MNA) was formed (ElBaradei Commission).The group was assigned to (1) identify and provide an analysis of issues and options relevant to multilateral approaches to the frontend and backend of the nuclear fuel cycle, (2) provide an overview of policy, legal, security, economic, institutional and technological incentives and disincentives for cooperation in multinational arrangements, and (3) provide a brief review of the historical and current experiences and analysis relating to multinational fuel cycle arrangements.In the report, MNA was assessed based on two primary factors, namely, assurance of supply and services, and assurance of nuclear non-proliferation.Furthermore, 3 potential MNA options were presented.
i.
To strengthen existing market mechanisms case by case with assistance from governments through long-term and transparent arrangement; ii.
To establish an international supply assurance such as fuel bank in collaboration with the IAEA as an organization to assure fuel supply; and iii.
To promote voluntary transformation of existing facilities of member states to MNA (including regional MNA by collaborative ownership and collaborative administration) The study results by the expert group at the IAEA are summarized in INFCIRC/640, which give an impact on the successive examination of multinational approach framework.After this report, a number of proposals related to supply assurance and multilateral approaches were put forward.The following are some of these proposals/approaches [15]: 1.In order to achieve "Reliable Fuel Supply (RFS) Initiative", announced by former Secretary of the US Department of Energy (DOE), Bodman in September 2005, the US is in the process of down-blending about 17.4 tons of HEU to about 290 tons of LEU (4.9%) within 3 years and storing them.The RFS Initiative was later renamed to the American Assured Fuel Supply (AFS) and it will be operational in 2012.
2.
During the discussion of fuel supply assurance at the Global Nuclear Energy Partnership (GNEP), the US, in collaboration with the partner states, declared that it would aim at establishing a fuel service mechanism including fuel supply at frontend and SF disposal at backend to achieve international nuclear non-proliferation.In the Nonproliferation Impact Assessment (NPIA) presented by DOE in January 2009, the importance of maintaining advanced reprocessing capacity including minor actinide recycling was insisted.It also emphasized the significance of the US's participation in the overall fuel services including backend service in order to suppress incentives for the emerging states to individually develop enrichment and reprocessing technologies.Later, being influenced by political regime change, the GNEP terminated its domestic activities (i.e.cancellation of prompt construction of commercial reprocessing facility and fast reactor) and decided that they would maintain international collaboration framework as International Framework for Nuclear Energy Cooperation (IF-NEC) only for international activities from 2010.The fuel supply working group at IFNEC expressed its willingness to support collaborative actions among member states and organizations towards establishment of an international fuel supply framework.It would also provide trustworthy and worth-the-cost fuel services/supply to the global market and provide options relating to the development of nuclear energy usage in accordance with reductions of nuclear proliferation risks.The new director expressed his speech its willingness to achieve so-called "from cradle to graveyard".
3.
World Nuclear Association (WNA) proposed a three-level assurance mechanism: 1) basic supply assurance provided by the existing market, 2) collective guarantees by enrichment operators supported by relevant governmental and the IAEA commitments, and 3) government stocks of enriched uranium product.According to them, it is necessary to promote the idea of an international reprocessing recycling center when nuclear energy usage is expanded in the future.
4.
Reliable Access to Nuclear Fuel (RANF) (nuclear fuel supply assurance concept by 6 states): Similar to the above, this proposal contains a three-level mechanism: 1) supply through market, 2) system in which enrichment operators would substitute for each other based on the collaboration with the IAEA, and 3) virtual or physical low-enriched uranium banks by a state or the IAEA.
Japanese proposal:
The states willing to participate shall voluntarily register at/notify the IAEA of their capacities (current stockpiles and supply capacity), and the member states shall notify the IAEA of their service provision capacity in accordance with the availability of service utilization capability by three levels (Level 1: provision of service on the domestic commercial basis -no exporting on a commercial scale, Level 2: international provision on a commercial basis, Level 3: storage that can be exported for a short time).The IAEA would make an agreement of standby-arrangement with member states and manage the system.If the fuel supply actually gets confused in a state, IAEA will play a role as a mediator.This proposal is to improve market transparency, prevent supply termination, and augment the RANF proposal.
6. UK Enrichment Bond proposal: Enrichment tasks shall be carried out by domestic enrichment operators.The supplying state, the consuming state and the IAEA will make a treaty in advance.The IAEA shall approve commitment of the consuming state for nuclear non-proliferation.If assurance is activated by bonding, the supplying state would not be prevented from supplying enrichment services to a consuming state.This proposal is to enhance credibility of supply assurance mechanisms and augment the RANF proposal.The Bond proposal was later renamed the Nuclear Fuel Assurance (NAF) proposal and was approved by the IAEA Board of Governors in March, 2011.
7.
The Nuclear Threat Initiative (NTI) proposal [16]: This is a storage system for LEU stockpiles possessed and controlled by the IAEA, and it is the anchor proposal for actual realization.For the activity of the NTI, the US pledged $50 million, Norway $5 million, the United Arab Emirates $10 million, the EU $32 million, and Kuwait offered $10 million.The total pledge has reached $107 million.Furthermore, in April 2009, Kazakhstan's President Nazarbayev announced that the country was ready to receive the IAEA nuclear fuel bank and officially announced its willingness to be a host state in January 2010 (INFCIRC/782).In May 2009, the IAEA presented a proposal for deliberation at the Board of Governors meeting held in June 2009.The proposal included consuming state's requirements in relation to the IAEA nuclear fuel bank, supply processes, contents of model agreement (e.g.supply price of LEU, safeguards, nuclear material protection, nuclear liability), etc.Later, at a regular Board of Governors meeting on December 3, 2010, the establishment of "nuclear fuel bank" which will internationally manage and supply LEU to be used as fuel for nuclear energy generation was agreed on.If the IAEA receives a request from a state which cannot purchase LEU due to exceptional circumstances impacting availability and/or transfer and is unable to secure LEU from the commercial market, state-to-state arrangements, or by any other such means, the IAEA will supply LEU to the state at the market price under the guidance of the Director General of IAEA.Through this agreement, the first system in which LEU would be controlled by an international organization began.The IAEA owns the bank based on the contributions from the member states.The Board of Directors will later deliberate the location of the bank.Kazakhstan is already declaring its candidacy to be a host state.
The resolution was proposed collaboratively by over 10 states including the US, Japan and Russia and was adopted with 28 states voting in favor.The developing countries which were planning to have nuclear energy later had been insisting that the bank would lead to a monopoly of nuclear technology by developed countries and "right for peaceful use of nuclear energy" stipulated by the NPT would be threatened.To address this issue, the resolution clearly stated that it would not "ask for abandoning" nuclear technology development by each state and obtained understanding from the developing countries.
8.
International Uranium Enrichment Center (IUEC) [17]: The IUEC was established in Angarsk, Russia, with investment by Russia and Kazakhstan.The IUEC is not only to assure supply but to provide uranium enrichment services.Thus, this proposal is more realistic than the others.The proposal states that the uranium enrichment technology will be black-boxed, namely, the investing states will not be informed, and the technology will be under the control of the IAEA.Other than Russia and Kazakhstan, Armenia and Ukraine are now members of the IUEC, while Uzbekistan is expressing its intention of participation.It will have the LEU reserve of two 1000MWlevel cores.In May 2009, for the deliberation at the IAEA Board of Governors meeting held in June, Russia submitted the proposal including the summary of agreement for LEU storage between the IAEA and Russia and summary of agreement for the LEU supply between the IAEA and the consuming states.In November 2009, being led by Russia, the nuclear advanced states submitted a resolution to the IAEA Board of Governors in November.The resolution was to seek approval of two agreement plans: 1) agreement plan between the IAEA and Russia to establish the LEU reserve under Russian IUEC, and 2) a model agreement plan between the IAEA and the LEU recipient states concerning the LEU supply from the reserve.The resolution was approved by a majority.In March 2010, the IAEA's Director General, Amano, and Director General of Rosatom Nuclear Energy State Corporation, Kiriyenco, sign-ed on the agreement for the establishment of the LEU reserve under Russian IUEC, and the LEU storage was established in December, 2010.
Multinational Enrichment Sanctuary Project (MESP) (proposed by Germany):
This proposal is for the IAEA to manage enrichment plants and exportation on an extra-territorial basis in a host state.The SNT will be black-boxed.
10.The Science Academies of the US and Russia presented analysis and proposals for nuclear fuel assurance as a measure to prevent proliferation of nuclear weapons under the title of "Internationalization of Nuclear Fuel Cycle -Goals, Strategies, and Challenges"13.In its report, the options and technological issues for the future international nuclear fuel cycle are presented.The report also contains the analysis of the incentives for the states that opt for accepting fuel supply assurance and developing enrichment or reprocessing facilities and do not opt for it.Furthermore, they examined new technologies for reprocessing/recycling and new reactors and made various proposals to the governments of the US and Russia and other nuclear supplier states to stop proliferation of SNTs and contribute to reduction in the risk of nuclear weapons proliferation.The report analyzed and summarized critical issues and presented several standards for assessing the options. .summarizes the flow of nuclear non-proliferation measures centered on multilateral approach/supply assurance in the past.As shown, the debates have become more and more active in recent years, and the need for internationalization of fuel cycle, which was not very realistic until now, is gradually becoming a reality.As described above, as of December 2011, the IAEA nuclear fuel bank, LEU reserve in Angarsk, Russia, and the UK's NFA proposal were approved by the IAEA Board of Governors, and the US's AFS begin its operation in 2012.
Issues with the past and current proposals
Most of the past proposed MNAs had never been implemented in any form until the nuclear fuel bank 7) and the LEU-IUEC storage 8) were approved by the IAEA Board of Governors.
It was probably because nuclear proliferation was not recognized as a sufficiently serious issue and there was not a very strong economic motivation.Many proposals included unfair double standards, i.e., "have" and "have not", and inconsistency with Market Mechanism.Also need of MNA may not have matured, or become critical yet.
However, as explained above, the situation has been changing in the last few years.Despite the Fukushima Nuclear Power Plant accident as well as the actual global concern over nuclear non-proliferation, the expansion of peaceful uses of nuclear energy in the world is unavoidable in the a long term.In that sense, the role of supply assurance was reviewed, and some of the above mentioned proposals have been approved by the IAEA Board of Governors.
Significance of MNA
Significance of NMA, namely, MNA's benefits and incentives of individual stakeholders may be summarized as follows: New nuclear non-proliferation regimes based upon mutual confidence and transparency, including regional safeguards, can be established, which can strengthen the function of nuclear non-proliferation.
Formulation of no discriminatory framework can be the primary incentive to make states join MNA.Recent criteria-based approaches of export of sensitive technologies in NSG [18] would help create a framework taking into account NPT Article IV.
Nevertheless, the number of enrichment and reprocessing facilities can be limited from the viewpoints of their needs (capacities) and nuclear nonproliferation, although every participating country can formally have the right to possess such SNTs.
Services on spent nuclear fuels, take-back, take-away, storage, reprocessing etc, can systematically be assured.Recipient countries can enjoy such services in NMA framework.
It is also expected that the host country in MNA would be discouraged to divert nuclear materials and to misuse related technologies because of the multilateral control of the fuel cycle facilities.
To minimize proliferation risks on SFs: accumulation of SFs, e.g., in power reactor user countries, has become serious issue in the world.By leaving such spent fuel in individual countries, there is also a certain level of risk to make such countries change the policy, i.e., to have an incentive to try reprocessing.
Improvement in 2S (safety/security) can be expected if for NMA framework systems to deal with such issues within the framework can be included, e.g., application of international standards among the participating countries.
Host countries may be able to expand their nuclear fuel cycle business capabilities further although facilities are expected to be controlled under/by MNA.
Prerequisites/features for establishing MNA
INFCIRC/640 (Pellaud Report) [19] proposed 7 elements of assessment called "Label" as prerequisites/features.In INFCIEC/640a variety of different issues are included altogether and the importance of each issue, such as nuclear security and safety, and political and public acceptance, is not focused individually, even though these are contemporary topics particularly following the Fukushima nuclear accident.Therefore, the following 12 elements, namely with 5 additional Labels, can individually be described as a full set of prerequisites or features to be considered for the formulation of a new framework of MNA as discussed elsewhere [20].
Label a: Nuclear non-proliferation
This includes safeguards, nuclear security and export control.If a state meets certain criteria (e.g.regional safeguards under MNA, nuclear security, export control), it is considered that the state can adequately maintain nuclear non-proliferation resume.Thus, the possession of sensitive nuclear technologies (SNTs) (i.e.uranium enrichment and spent fuel reprocessing), which is one of the measures for nuclear non-proliferation, would not necessarily be limited (criteria-based approach).
Label b: Fuel cycle service
An appropriate state becomes a host state or siting state (that provides/lets site) and offers fuel cycle service, based on the above-mentioned criteria-based approach.It includes uranium fuel supply service and services on spent fuel treatment.The latter should be made with a clear plan/agreement for long-term spent fuel treatment (storage / recycling / direct disposal), e.g., reduction of nuclear waste toxicity (HLW to medium level), individual member states to receive final waste, and use of MOX, in order not to bring concern to host/siting states.
Label c : Selection of a host state (siting state) The state that meets all the criteria can be a host state or siting state.The specific criteria to participate in the multinational framework or to be host/siting state are, for instance, to sat-isfy conditions almost equivalent to the "objective criteria" described in INFCIRC 254 part 1-6, 7 (NSG guideline revised) 18, that is, member states are in full compliance with its obligations under NPT/safeguards agreement, are adhering to the NSG Guidelines, apply agreed standards of physical protection and have committed to IAEA safety standards.
Label g : Transport
Member states should cooperate and maintain international standards for nuclear material transportation beyond borders.
Label h : Safety
International safety standards should be met within MNA.
Label i : Liability
MNA should cover a certain level of liability.
Label j : Political and public acceptance
Individual host or siting states should obtain political and public acceptance in corporation with MNA.
Label k : Geopolitics
Practically it should be taken into account if the stat is geopolitically stable.
Label l ; Legal aspect Table 1 summarizes the existing treaties and agreements that correspond to each Label to be considered for MNA.
The gap between the new-MNA and existing related laws and agreements, which may conflict in some cases, should be adjusted.In particular, new MNA framework must have equal or higher capability on nuclear non-proliferation (Label a), e.g., in order to adjust the existing bilateral agreement that may be one of the strongest measures among the existing nonproliferation systems.In other words, the MNA member states must basically assure conditions set forth in the international treaties and agreements.
An example of specific MNA framework study [21, 22, 23, 24, 25]
The author's group has been studying an example of MNA framework, where strengthening of international non-proliferation scheme and provision of stable energy/nuclear fuel cycle services in a region are discussed.It contributes to enhancement of transparency and trustbuilding in the region.The study investigated the schematic issues and the countermeasures concerning the specific measures to achieve stable maintenance of the multilateral international nuclear fuel cycle including stable uranium supply system, spent fuel treatment system, usage of plutonium, establishment of regional safeguards scheme for the international nuclear fuel cycle, requirements for an organization that carries out international nuclear fuel cycle, and roles of industry in the international nuclear fuel cycle scheme.An image of framework scope is given in Fig. 4. Outline of the study is shown below, Three options on MNA system, Type A, B and C as shown below are defined.
Type A: No involvement of services (assured) of fuel supply, spent fuel storage and reprocessing, but regional framework for 3S.
Type B: Provision of services (assured) of fuel supply, SF storage and reprocessing without transfer of ownership of facilities; including regional framework for 3S.
Type C: Provision of services (assured) of fuel supply, spent fuel storage and reprocessing, MOX storage with ownership transfer of facilities to MNA; regional framework for 3S (with IAEA -arrangement).Regarding SFs, the MNA consisting of host, siting, and recipient states has clear plan for long-term spent fuel treatment, i.e., recycling / direct disposal, reduction of nuclear waste toxicity and use of MOX, within a specific certain period in order not to bring concern to host/siting states.
MNA develops technologies and services of reprocessing to reduce radio-toxicity of HLW (e.g., HLW to medium level) that would make an individual receive final disposal waste easier.It establishes Regional Material Accounting and Safeguards system within the MNA Framework to implement the nuclear non-proliferation regime, as described in Fig 5 .The MNA agreement contains high level of nuclear non-proliferation capability, equivalent to the existing bilateral agreements (e.g. one with the United States).
The MNA has function to attain the international level on nuclear safety and security for facilities within the Framework (not only for fuel cycle facilities but also nuclear power reactors); criteria and inspection system.
The MNA has agreement with technology holders to precisely manage and control the SNTs (limited to technology-holding operators only).
Obligation with regards to nuclear non-proliferation is performed equally by the member states, while it is guaranteed that the right of peaceful uses of nuclear energy pursuant to Article 4 of NPT is not interfered with.
The specific requirement to participate in the multilateral framework is to satisfy conditions almost equivalent to the "objective criteria" described in INFCIRC 254 part 1, 6-7 (NSG Guidelines revised in 2011, see below*).
The MNA Framework is to be more economically advantageous than the fuel service on per state basis.
Framework member states cooperate on and agree to "transport" with regards to the nuclear fuel cycle service.and geopolitical perspectives)", "to realize international standards for safety and security", "to have higher economic potential for fuel cycle than a single state approach", "to eliminate conflicts/inconsistency with existing laws and regulations", and "to solve transport issues of nuclear fuel, etc".Particularly, involvement of industry would be a key issue when such incentive or attractiveness of the new proposal is discussed.Taking into account the example of EURATOM, the need of MNA for participants is the overriding issue in having "incentive" towards establishment of MNA.The author would like to note that the environment is getting ripe for the need of MNA, in terms of SF and waste treatment, maintenance and improvement in safety, security and nuclear non-proliferation-safeguards (3S).
There are still many challenges, pursuing incentives on economic efficiency, 3S, finding the solutions for nuclear material transportation within MNA, effective and efficient organizational management with not only member states but industries, and legal conflicts between new MNA's treaty/agreements and existing ones.
Conclusion
Even after Fu-NPP accident, use of nuclear reactors may be expanded particularly in emerging countries, where reliable systems of fresh fuel supply as well as proper management of spent fuels, free of any political disruptions to their nuclear reactors, are highly desirable.Establishment of international cooperative systems, which includes services for fresh fuel supply, spent fuel take-back/take-away, interim storage, reprocessing, and possibly repository disposal, may be able to contribute to a) enhancement of 3S, nuclear non-proliferation (Safeguards), Safety, and Security, b) economic rationality, c) promotion of confidencebuilding, and d) prevention to the occurrence of unfair business such as government-to-government transaction based on cradle-to-grave service that nuclear weapon state's privilege enables.This kind of internationally cooperative framework may become essential for future sustainable utilization of nuclear power.
Fig. 2
Fig. 2 Diversification of Uranium Resource Supply
Figure 3 .
Figure 3. Transition of Proposals/Initiatives for International/Regional Management of Nuclear Fuel Cycle Relevant to Nuclear Non-proliferation
Figure 3
Figure3.summarizes the flow of nuclear non-proliferation measures centered on multilateral approach/supply assurance in the past.As shown, the debates have become more and more active in recent years, and the need for internationalization of fuel cycle, which was not very realistic until now, is gradually becoming a reality.As described above, as of December 2011, the IAEA nuclear fuel bank, LEU reserve in Angarsk, Russia, and the UK's NFA proposal were approved by the IAEA Board of Governors, and the US's AFS begin its operation in 2012.
Label d : Access to technology Particularly, the access to SNTs should be strictly controlled under the MNA Framework.Label e : Multilateral involvement This includes 1) having multilateral cooperative system on e.g., safeguards, safety, and security, and 2) provision of services with or without transfer of facility ownership to MNA.Label f : Economics MNA, as a whole, increases economy when compared with management by individual states.
Fig. 4 .Figure 4 .
Fig. 4. Possible Framework of Future Nuclear Fuel Cycle Figure 4. Possible framework of Future Nuclear Fuel Cycle
(
Specific framework proposed) It is targeted to the Asian region It establishes MNA Operating Organization as the core of the framework function Conclude Treaty on Regional NFC and related Agreements between States and the Organization.In the multilateral framework, the system/facilities are divided into the Type A, Type B and Type C. Plutonium-handling facilities such as reprocessing, MOX fuel fabrication, Fast Reactors, and MOX storage facilities should be controlled under type C, whereas uranium enrichment facility and spent fuel storage can probably be categorized as Type B or C depending on the siting countries.LWR MOX reactor would be Type A, while direct disposal should be Type B.
Probably, nuclear societies including industry have internationally received greater recognition of the importance in Safety, Health and Radiation Protection; Physical Security; Environmental Protection and Handling of Spent Fuels and Wastes; Compensation for Nuclear Damage; Nuclear Non-Proliferation and Safeguards and Ethics, as described by Principle of Conduct [28], since the Fukushima Power Plant Accident. | 9,317 | 2013-02-06T00:00:00.000 | [
"Environmental Science",
"Engineering",
"Political Science",
"Physics"
] |
Islam and Prejudice: Special Reference to Gordon W. Allport’s Contact Hypothesis
This study explores the Muslim perspective on human interaction, relationships and prejudice. A survey of the literature recognises Islam’s fundamental acknowledgement of human diversity, drawing on a dynamic theological, moral, spiritual and legal philosophy revolving around the preservation and sustainment of non-prejudiced human contact. This study discusses the Muslim perspective of human contact, non-prejudice, and accordingly, revisits Gordon W. Allport’s “Theory of Contact Hypothesis” in an effort to compare and contrast it with the Muslim perspective on related issues such as racial prejudice, gender inequity, age prejudice, disability discrimination, social status and classism. This research concludes that Islam has developed a framework necessary for cultivating religiosity and morality without risking the value of effective and harmonious human relations. Further empirical studies on the interplay between Muslim theory and practice on contact hypothesis and prejudice are required to further interpret the dynamics of Muslim values in working settings, and the viability of translating religious ideals into reality.
Introduction
Prejudice can be described as the unfair and rigid judgment about others. According to Woolfolk, prejudice is constituted of beliefs, emotions and tendencies towards particular actions (2016,251). This cognition initiates from childhood in the form of schemas, evolving as stereotypes in such patterns of belief against a group of people throughout a lifespan. Research on the connection between prejudice and religion began with early studies of Adoree and his colleagues in the late 1940s (Hood Jr., Hill and Spilka 2009) and continues to support similar findings (Rowatt, Carpenter and Haggard 2014). The conceptualisation and measurement of religious orientation however, began with Allport and Ross (1967) and continued with Gorsuch and McPherson (1989) in the context of prejudice. Allport is recognised as the founder of the cognitive approach of prejudice (Dovidio, Glick and Rudman 2005) over discussions on the nature of prejudice in the field of religious fundamentalism. Allport views prejudice as a negative feeling or attitude, a failure of rationality (Allport 1966, 448) and a stereotyped overgeneralisation (Allport and Ross 1967, 412). In a religious context of prejudice, Allport sought to investigate the varying degrees of prejudiced belief of church attendants and nonattendants, both on the theoretical and empirical level. According to him, church non-attendants exhibit less prejudice than church attendees. Allport (1966) believes that communal and extrinsic religion become particularly imperative in the context of theological prejudice. For him, religious orientation of extrinsic and intrinsic motivation which represents a "subjective formation within personal life" (Allport 1966, 456), and sheds light on the varying degrees of prejudice among church groups; though the extrinsic and intrinsic conceptualisation was revised to adapt a better understanding of religion (Gorsuch and McPherson 1989). For instance, church attendance was perceived as an intrinsic rather than extrinsic orientation. Interestingly, the findings of both Allport and Ross (1967) indicate that the intrinsically religious are less prejudiced than their extrinsic counterparts, 1 implying that the extrinsic orientation draws on a strong support of intolerance and bigotry whereas the intrinsic orientation favours tolerance. Moreover, an extrinsic religious orientation is found to be compatible with prejudice whereas the intrinsic orientation supports tolerance and humanitarianism. More particularly, the extrinsic orientation focuses on the outlook of religion rather than perceiving the meaning of religion per se and hence is perceived as a weak superficial act, subject to easily changed belief (Allport 1966).
Five decades following Allport's research, the theological context of prejudice remains intriguing whereas interesting findings remain beyond the idea of communal and extrinsic religion seen as a cause of bigotry. In their work, The Psychology of Religion: An Empirical Approach, Hood Jr., Hill and Spilka (2009) presented a broad perspective of religion and prejudice according to which religious orientation should be reflected through a comprehensive meaning of culture, society and other relevant factors. Interaction and contact with other cultural and social systems need to be examined in an objective evaluation of the relationship of religion vis-à-vis prejudice (Ghorbani et al. 2002). Since major religious traditions support or motivate their adherents to love one another, thus resulting in reduced prejudice, religion again plays an important role in the mechanism of orienting systems of perception and engagement with others at large (Randolph-Seng 2014).
Allport's proposal on contact hypothesis is regarded as instrumental to understanding the religion-prejudice relationship. According to the contact hypothesis, positive attitudes increase when contact between group's increases, while prejudice may reduce among opposite groups through out-group contact under optimal conditions (Allport 1954). This hypothesis demonstrates that the understanding of religious tradition is critical not only to the clarification of the role of group contact, but also in regards to the position of religious beliefs, spirituality, morality and laws concerning the philosophy and quality of intergroup contacts. As such, this study seeks to explore the Islamic perspective on group contact with special reference to Allport's hypothesis in view of the assumption that religious approaches to prejudice may offer insightful understanding leading to an enriched understanding of human contact and inter-group interactions. More specifically, the exploration of the Islamic insight on intergroup contacts is useful in view of the rising degrees of tension, conflict and prevalent marginalisation of religious communities in the present day. The Islamic perspective remains significant, with respect to not only understanding the other, but also in view of incorporating diverse global religious and native perspectives, and capitalising on Islamic values and principles towards engineering realistic, effective and more inclusive approaches to the study of human intergroup contact. This study seeks to explore the relevant religious Muslim texts to gain an understanding of the normative religious notions that define Muslims' perception of others and how adherence to the Islamic religion shapes its followers' attitudes towards prejudice, or threatens out-groups.
Islamic Perspective on Human Contact
This section discusses the Islamic perspective on human contact and highlights the underlying factors affecting Muslims' interpersonal and intrapersonal perceptions, relations and rapprochement. The understanding of the theological and religious premises vis-à-vis the human contact theory is critical to appreciating the nature of Muslims' engagement in cross-cultural and inter-communities' interactions, while detecting concepts that may infiltrate the web of Muslim beliefs and thoughts. The methodology in this study is primarily one of textual analysis and hence normative. First, one finds that the Qur'anic view on human relations is imbued with the theological foundation of the unity of creation and the principle that humans are children of Adam, and that none are entitled to privileged treatment over others. As such, there is no effectual value assigned to human characteristics such as race, colour, gender, age, language, social status, physical appearances and the like. It is with this view that the Qur'an exhorts its followers to initiate spaces for dialogue, interaction and courtesy. 2 Muslim sources discuss unity, sameness of human race, and further underscore the imperative for activating spirituality and religiosity to preserve the moral quality compass. Moreover, while they hold belief in human sameness, they appear to resist all perceptions of assigned human characteristics or attributes, maneuvered racial or socio-economic differences, and instead acknowledge good contributions to life alongside a profound appreciation of qualities of knowledge, character and righteousness. Sameness therefore implies a reference to the original creation, reflecting the equality of human species before God, and the rejection of all forms of distortive representations of human contact such as humiliation, exclusion, abuse, exploitation, control, manipulation, sacrifice or exaltation.
Among the renowned Islamic religious directives in this regards is a prophetic declaration that people are as equal as the teeth of a comb, bringing to mind the image of a human family line tracing back to the progenitor Adam, and further cautioning against prejudice, and racism which carries no self-merit. In addition to this concept is the principle that difference of skin colour holds no social implications, and that the pride of early Arabs in the pre-Islamic period towards clan and lineage is outright rejected. Many similar religious directives require Muslims to maintain courteous interactions with their co-religionists, give in charity with sincerity, treat neighbours with kindness, conduct business with trust and integrity, lend and borrow with justice, choose spouses without racial bias, treat employees, clients and customers justly and decently, avoid ridiculing and nicknaming or mocking others, and set equal social justice fairly (Qur'an 5:42).
Islam also revisits the concept of beauty and humanity, making it worthwhile to draw attention to its inner dimension, and setting the perception, evaluation, interaction and norms of human treatment accordingly. Such philosophy places less attention on colour, race, language, wealth, gender, physical strength, ancestry, traditions; and instead places emphasis on the universal criterion of character, piety and goodness. Drawing on the Qur'an's perspective, one finds use of the term ta'aruf (knowing one another) which is derivative of ma'rifah (knowledge), and this places focused attention on the level of character building, cultivation of piety and self-discipline; all of which highlight Islam's prioritised interest of the inner human character and the intrinsic beauty of man as illustrated in numerous traditions of the Prophet Muhammad (pbuh). 3 Setting human judgement in a rather vertical axis inclined towards divine satisfaction is perhaps purposeful towards molding human attitudes away from superficial or artificial considerations, and instead, around inner persuasions with effectively objective effects on their surroundings.
This theological perspective adds further value to the quality of human contact by way of drawing differences to their original circles, all without neglecting core humanistic values, and further reminds of the importance and often-neglected innate human dimension resulting from inattention to the spirituality of man. Interestingly, the Qur'an revisits the filters of race, colour and ethnicities, addressing them in light of the principles of unity of creation and the sameness of creation, with the purpose of eliminating accumulated judgments, biased assumptions, misperceptions and returning to the beginning of human contact. To this end, the tradition of Prophet Muhammad states: All of you descend from Adam, and Adam was made of earth. There is no superiority for an Arab over a non-Arab nor for a non-Arab over an Arab, neither for a white man over a black man nor a black man over a white man. (Albani 1984, 361) Along similar lines of thought, Islam is keen on developing Muslims' knowledge and familiarity of other communities and understanding their conditions, backgrounds and languages. This is likely set to nurture mutual acquaintances, drawing groups closer and eliminate human hazards. It is with such a view that Islam seeks to enrich intercommunity experiences and tolerance, promoting knowledge of the out-group, boosting community contact and further increasing empathy and perspective taking.
Prior to forging intercommunity contact, Islam seeks first to shape the ethical and spiritual in-group contacts, and as such mandates a basis of moral rectitude and established religious congregation as an effective arrangement for rapprochement so as to strengthen intergroups' contact while enhancing the ethical attributes of interconnectivity, integrity, tolerance and courtesy. As for intergroup contact in Islam, there exists a series of religious recommendations to consolidate the community's spiritual and social capital, and further enhance their effectual contact while increasing contact frequency and building on the resulting mutual discovery. This is represented for instance in the Muslim daily congregational prayers; mandatory Friday prayer; Eid prayers; Ramadan nightly prayers; funeral prayer; prayers of eclipse, rain, pressing need and so forth. In addition to the Five Articles of Islam, Islam also sets its ethical code of hygiene and public health, diet, business ethics, family management ethics, social solidarity, human relations and others to sustain an effective capital driving the vision of Muslim communities and ultimately their intergroup contacts.
The description discussed above should not suggest any ideal thesis of Muslim intergroup contact but rather anticipates humans to act upon their free will whilst exhibiting ordinary behaviour of compromise, conflict or reconciliation. This may also give reason as to why even in the presence of religious and spiritual discipline, human conditioning by way of education and law is critical in ensuring non-prejudiced intergroup contacts. In view of the above, however, there is a need to underscore Islam's keen interest in forging broader objectives for intergroup contact and exchange with a philosophy inherently imbued with the humanity, belonging to God and shared belief throughout the Adamic creation.
Alongside the significant interest in intergroup contact and fulfilling of groups' fundamental needs, Islam also introduces relevant intergroup principles best identified as the common ground (Qur'an 3:64), which in the context of religious intergroup contact, seeks to nurture interactive contacts leading to effective mutual familiarity, understanding and healthy intergroupness. The interest in common ground however, may turn disparaging with active prejudiced intergroup contact (Qur'an 30:22) in view of the fact that the commonality mindset opts for naturalised relations, compromising attitudes, mutual recognition and inclusive participation. Common ground engagement affects the intergroup contact, not necessarily via dissolving one's own religious beliefs or rituals, but rather, through the will to adopt differences, celebrate diversity, capitalise on shared human values of dignity, equality and humility, collective negotiation of the religious space, and deepening consciousness in regards to a "civilised" cross-religious unity. At the most generic level however, this should translate into a serious interest in education, interaction, communication, consolidation of power against humans, and social and environmental problems while constantly transforming spiritual inputs into constants of goodness and pity.
In addition, and congruent with the Islamic logic of human contact, one notes that in Islam there surges a profound yet compelling exhortation to assist others indiscriminately and unconditionally, while continuing to validate one's religiosity and spirituality by way of service, assistance and sympathy. The Qur'an is replete with exhortations to such meanings and virtues, as in the case with verse 5:2 which exhorts believers to cooperate on righteousness, and deeds that are beloved to God, whether related to individuals or groups. Religiosity is thus placed in the framework of freeing the self from egoistic self-centeredness and the actively passionate giving to others. However, what appears to be critical in the context of discussion at hand is that such a religious perspective also implies borderless resourcefulness, conditioned neither by belief nor by race, but rather applied to all openly. It is with this philosophy that Islam stresses not only furthering intergroups' contact but more importantly raising the quality of the contact itself. Resourcefulness transcends the colourful lenses of ideological, cultural or socioeconomic assumptions. Islam appreciates other groups, recognises their needs, and hence views marginalisation, human indifference and exploitation as a grave fault. To further protect the spiritual, moral and social capital of the community, it sets preservation of belief, life, intellect, human progeny, property, honour and dignity, environment, justice, freedom and social-human relations as a higher ends for good life. Those Islamic ideals represent the kernels of community building which bonds members to common living principles as they seek to nurture and sustain common interests in belief, health, life, intellect, wealth, freedom, justice, environment, and safeguard effective intergroup contacts, resulting in sustainable levels of interaction and exchange.
According to the Islamic viewpoint, these ideals are reflective of the innate human predisposition, intellect, human experience and revelation which seeks to ease life, sustain comfort and to avoid pain, escape inconvenience and harm, and add value to life, health, wealth, intellect, human lineage, freedom, justice, human relations and the environment, while improving and safeguarding those ideals and sustaining happiness. In this regard, rationality of religion posits that for sustainable living, people need to uncheck all barriers and prejudices that may jeopardise their shared interests, and set laws for life that ensure human interconnectivity. The legislation of the common interest is advantageous not only to the well-being of the community, but also to capitalising on human needs, aspirations and the desire to enjoy decent degrees of dignity and justice.
Islam and Non-Prejudice
In light of its integrative nature, Islam lays out an ethical and social order, resonating with innate human predisposition, intellect and community aesthetic. This, for instance includes the obligation of extending a helping hand to the less fortunate, giving in charity, caring for one's neighbours, 4 mentoring and guiding others, sharing sustenance and open door homes, 5 engaging in business with others, lending money, showing leniency, acting with trust and forgiveness, visiting the sick, attending funerals, and avoiding abuse and exploitation of others. Islam is also concerned with developing amiable interpersonal interactions, and as such, dictates rules prohibiting mockery, nicknaming, backbiting, gossiping, looking down upon others, thinking ill of others, and foregoing malicious, divisive and harmful speech. 6 Similarly, it exhorts its adherents to demonstrate integrity, nobility, trust, confidentiality, decency and justice, even in cases of extreme enmity, 7 to avoid disrespectful argumentation, or to insult others' faiths. Moreover, Muslim sources confirm a strong commitment to the community aesthetic, taste and the need for an atmosphere conducive to the enhancement, comfort and consolidation of intergroups contact. Chief among those expectations are hygiene, cleanliness, green conservation of environment, food diet, personal care, cheerfulness, adopting positive names, avoiding superstition and pessimism, 8 personal grooming, 9 public etiquette, animals care, avoiding witchcraft and magic, speaking gently, and walking humbly and decently. In the following section, we shall explore the objectives of Islam set to reduce prejudice in the context of the contact hypothesis.
First, the study of the principles of Islam points to a large interest in placing differences in the purview of diversity (Qur'an 109:6, 30:22), with an encouragement of mutual acquaintances (Qur'an 49:13) and of human equality, alongside exhortation on cooperation (Qur'an 2:5) and kindness for all (Qur'an 2:83). In Islam, human physical characteristics may constitute the person it qualifies as a unique being, but is nonetheless irrelative to the person's creaturely status before God. Such characteristics never determine ethical outcomes and are never conclusive, for it cannot be an axiom that a person with any imaginable combination of such characteristics is morally worthy or unworthy. The inner core is constitutive of the person's being and must therefore remain, to the extent required, free of those characteristics; both capable of using their determinative power as well as of doing otherwise, i.e., of channeling their causal efficacy to other ends (Abusulayman 1989, 46-47). As such, Islam holds that righteousness and its associates of goodness define one's self-image and value, thus rendering social status reflective of virtue, refined character and good contributions to others as shown in many religious traditions. 10 Islam also seeks to maintain equilibrium in human life while acknowledging human needs, passions and aspirations, which inevitably result in difference and distinction.
Interestingly, Islam does not appear to be in favour of disregarding human difference or demeaning people's earned merits, but instead preoccupied with disciplining the self and cherishing virtues of thankfulness, justice and compassion. For instance, Islam holds a high regard for work, ownership and earning, yet not according to an egalitarian basis of economy, ownership, living standards and merit. This leads to resistance against all forms of exploitation and degradation, preventing classification of recognition and merit according to gender, physical characteristics and other traits. Similarly, Islam censures the dissension of social classes and social conflicts in seeking social justice; but rather endeavours to build harmony through equality and justice as shown for example in its resistance against the monopoly of wealth and resources.
Second, the Islamic definition of gender falls within its broader Weltanschauung on the human creation and sameness, as seen in the definition of Muslim Revelation delineating the position, duties and responsibilities of both man and woman. The overwhelming majority of Islamic literature views man or women not according to their sexes and gender, but rather as the individual cumulative of humanness and accomplishments. According to Islam, gender alone does not grant special privileges or preference (Qur'an 4:19); piety and moral character however do, as shown in the following Qur'anic verse: "Whether male or female, whoever in faith does a good work for the sake of God will be granted a good life and rewarded with greater reward" (Qur'an 16:97, 33:35). The hadith tradition is replete with narrations on the principle of gender constructive partnerships, which is to be experienced with love, compassion and kindness (Qur'an 30:21). 11 Along similar vein of thought and drawing upon the fundamental normative ethics of Islam on human equality, dignity and honour, the prodigious literature expounds on Islam's value for old age in granting a close bearing with God, thus providing it with theological and spiritual content. 12 In his classic definition of ageism as "a process of systematic stereotyping of and discrimination against people because they are old" Robert Butler goes on to observe: "Old people are categorized as senile, rigid in thought and manner, old-fashioned in morality" (1975,11). Ageism is often defined as prejudice and discrimination against older people based on age and includes for example denial of resources or opportunities, or viewing old age in negative stereotypes. The Islamic traditions give prominence to honoring old age, which further alludes to a distinct spiritual character of the elderly by means of distinguishing sanguine descriptions of grey hair taking the form of radiant light on the Day of Resurrection (Tirmidhi n.d., 3:224), and a sign of ascendancy of religious stature (Baghawi 1991, 6:211).
Islam maintains that age itself is neither demeaning nor an infirmity resulting in maltreatment, discrimination or prejudice. Islam acknowledges several weighty religious and spiritual dimensions corresponding old age with close acquaintance to the divine, acquisition of wisdom and elevated honour. As such, old age is not viewed as a burden but rather an avenue to a rewarding religious experience and drawing close to the divine. In fact, Muslim culture, which grooms its youth to show respect, defers to older people and further treats them with honour and dignity; in turn generating high social regard and respect for older individuals. 13 Islam also establishes a number of spiritual, ethical and legal measures in ensuring just treatment and mercy showcasing essential values of dignity, honour, kindness, respect, appreciation, ease, support, solidarity and service (Bensaid and Grine 2014). All this is in addition to establishing justice over their rights and needs while guaranteeing an equitable distribution of responsibility in a manner that transcends religious, ethnic and cultural demarcations and discernment (Bensaid and Grine 2014).
In light of the above, one understands the spirit of Islam towards old age and elderly, and how it constitutes an ethical and legal system around dignity, respect, care and justice. Interestingly, those Islamic preventive norms transcend race, religions, classes and colour; alternatively revolving around the very essence of human creation, and as such, acknowledging inevitability of ageing along with its implicit occurring changes in physical appearance of the face, body, speech, mobility, cognitive abilities and emotional reactions. Moreover, Islam underscores the religious and moral duty of respect and care for elderly parents and the imperative to show them mercy and kindness in addition to safeguarding the financial rights of even non-parent elderly people. The sum and substance is that old age carries no demeaning association nor does it lower the image of the elderly in the community. Islam seeks to conjointly ensure a life of respect and dignity and to enrich the community's spiritual and moral culture high esteem for the elderly, in which they may oversee social functions, acting as witnesses to marriages, conflict, divorces, birth, social mediation or reconciliation.
Islam also advances its standpoint on people with disability which reflects profound theological persuasions reflective of the belief that true disability is that of the conscience and heart (Qur'an 22:46), and not of the body, in addition to viewing physical disability as rewardingly divine tests, 14 in need for discipline of patience and thankfulness. Disability inspires the community to demonstrate thankfulness to God over their health and good shape, to show respect (Qur'an 49:11), console people with disability, 15 socialise, share sustenance (Qur'an 24:61), accept their invitations, integrate them in society, 16 pray for them, avoid deluding them, not show contempt, appreciate their qualities, join them in marriage, ease their inconveniences as well as remove their hardships (Qur'an 48:17). Physical changes may be irreversible and should be welcomed, and reacted to with utmost sensitivity, care and appreciation. In its place, it is believed that spiritual health, intellectual soundness, and social maturity would render age within a paradigm of natural prospects of change, and bring about the most humane and rational treatment of.
Islam and Group-Norm Theory of Prejudice
Allport's contact hypothesis states that the decline of prejudice directly corresponds to the amount and type of contact occurring between differing groups or cultures. The type of contact that mediates prejudice is the most influential aspect for initiating attitude and behavioural change. Allport sought to explain how the religious context instigates prejudice. In this, there exist mainly three categorisations of religious prejudice: the theological context, the socio-cultural context and the personal-psychological context (Allport 1966). Among those three categorisations, Allport draws attention to the context of intrinsic and extrinsic religious orientation of "person-psychological context". Allport (1954) further contends that interaction may lessen prejudice: (1) When those in contact are of equal status, (2) when those in contact perceive themselves to be in pursuit of common goals, (3) when contact is sanctioned by institutional support, the larger community or law, and (4) when the contact is positive in nature. Such interracial contact is most possibly apparent in diverse cultural societies (Schmid and Hewstone 2010). This contact could be evaluated again with a more specifically personal commitment and individual differences based on Allport's religious orientation. It appears that the findings of intrinsically motivated people exhibiting less prejudice implies that living the meaning of religion may in large part be advantageous to promoting equal status and cooperation, leading to common goals and institutional support with the implication that intrinsic orientations contribute significantly to intergroup relations. In Islam, belief in God and doing good (e.g., Tekke et al. 2016) are imperative for an understanding of intrinsic orientation, which could be a process of the personal-psychological context, expounded into relationships between commitment and religion. Therefore, good deeds are found here to affect the accountability of religious people to society (Tekke et al. 2015).
This social accountability helps maintain relationships with others as shown in Islam's interest to promote cooperation and intergroup relations through exhibiting resourcefulness to other groups. Good deeds, however, are not necessarily prerequisites for positive interaction towards out-group members seeing that religious groups, especially fundamentalists, are interested in the welfare of their members alone and not of out-groups (Theiss-Morse and Hibbing 2005). Therefore, it is argued that a prevalence of individual differences and negative attitudes might lead to prejudice (Whitley and Kite 2010), for the reason that one of the most frequent sources of prejudice is rooted in needs and habits that reflect the influence of in-group memberships upon the development of the individual personality. A member of a religious group associated with "fundamentalist" values may pose threats to other belief systems. Such prejudice is greatly associated with fundamentalist views, according to Allport (1954, 9), who asserted: "[prejudice] is actively resistant to all evidence that would unseat it".
In light of good deeds and social accountability, Gaertner and Dovidio (1986) enumerates the following postulates of Allport which he sees as necessary to increase intergroup cultural openness: Interactions among members who do not possess qualities that are stereotypically associated with their group membership; situations that provide strong normative and institutional support for the contact; similarity of beliefs and values between groups; opportunity for intimate self-revelation, and personal contact. (318)(319) According to Brown (1995), social and institutional support are measures employed to promote greater contact and interactions, existing to help create a social climate ripe for tolerant norms to emerge. According to Pettigrew and Tropp (2006), institutional support is a vital condition leading to the reduction of stereotypes and prejudice. Moreover, the contact between groups must be long enough, often enough and close enough for proximity to affect attitudinal change about different groups (Brown 1995).
Yet to help promote equality and cooperation among different groups, belief in God, social accountability, institutional support, individual difference and intrinsic motivation ought to be considered when carrying out good deeds . Researchers like Batson and colleagues have revised Allport's conceptualisation of intrinsic orientation from belief and cognitive activities to social accountability and institutional support. The revised conceptualisation of religious orientation proposed by Gorsuch and McPherson (1989) factors in behaviour and religious attendance, often exploited by researchers. This understanding appears to be compatible with the Islamic perspective.
Conclusion
Despite the unwavering commitment of the Islam with its spiritual, moral-ethical, legal systems to cultivating and sustaining non-racial treatment of ethnicities, races and languages, and to eliminating all attitudes, practices and forms of ethnic prejudice and racism, the fact obviously points to continuing gap between theory and practice. This theoretical study does not intend to advance kind of utopic image of Muslim society free of all abuses, discriminations and racial practices. This is because Muslim society and culture continue to exhibit violations of those very norms and values, which in turn, could be caused by other complex factors including poor education, culturally bound religious textual interpretations, political maneuvering, sectarian interest manipulations, colonial exploitation, narrowed approach to change and so forth. One such present example is the notion underlying the privileged position of Arabic and Arabism, and in some instances, the related preference towards of Arabs over non-Arabs; this is despite the abundant textual evidences, which strongly call for abolishing all forms of linguistic preferences and rather insist on piety and righteousness as fundamental criterion.
Another related example is Islam's bold call for justice, kindness and courtesy towards other faith groups. The Islamic literature is replete with textual evidences underscoring the imperative to treat other faith groups with justice and fairness, and to act fairly, kindly and courtesy towards them. Dwelling on those texts would exceed the limits of this paper. Clearly, the theological teaching of Islam does not pose problem to the possible racial treatment or abuse of other faith groups. What appears to drive groups' behaviours is the problem of interpretation on the one part, and hence, there needs to be ongoing diagnosis of the forces underlying community dynamics. The religious and spiritual capital can certainly help remedy abnormal practices affecting human relations, however, to fully benefit from those capitals in the redefining and engineering of new web of relations and human interaction culture, one needs to re-evaluate the context and highlight the complex factors that directly and indirectly contribute to the shaping of current culture.
Nonetheless, through its integrated approach of the intrinsic and extrinsic values related to both in-group and out-group settings, Islam seeks to nurture individuals' and groups' resilience against prejudice and to sustain their immunity from discrimination. On the internal side, Islam deepens belief in individual accountability before God, accountability for one's acts of prejudice, and appreciation of dignity, equality, honorability of the human race. The translation of Muslim values and principles on sustained intergroups' non-prejudice is highly dependent on factors of education, engineered culture, methods and approaches of religious interpretation, styles of religious indoctrination and learned cultural stereotyping. Prejudice can be un-learned, and while spirituality goes hand in hand with learning, it invokes moral, ethical and legal commitment, in both the individual and public spheres to sustain the learning of non-prejudice. This process of un-learning prejudice and learning of non-prejudice can be characterised as profound, active, transparent and reflective in view of the very nature of the process and discipline of spirituality and religiosity in the life of the Muslim.
Nevertheless, it appears that the rediscovery of Islam's fundamental values of spirituality and religiosity could provide the framework necessary for shaping and at the same time nourishing much of Muslim's attitude, behaviour and feelings, and as such their interpersonal and intrapersonal relationships. This view continues to gain increased recognition in academia despite the ongoing problems of violence and terror allegedly associated to religion and to God. This current research has shed light on human contact, relationship and prejudice. Confirming Allport's contact hypothesis, Islam promotes such an environment of non-prejudiced to community to develop mutual understanding and co-existence. However, the suggestion of Muslim religious orientation of human contact, more particularly the intrinsic and extrinsic orientation, might be beneficial in spite of the present measure of religious orientation (Allport and Ross' religious orientation) being validated across a several countries. Muslim measures such as the "The Religiosity of Islam Scale" of Jana-Masri and Priester (2007) and "Hoge Intrinsic Religiosity Scale in Muslims" by Hafizi, Koenig and Khalifa (2015) help evaluate spirituality and religiosity, unfortunately they neglect the broad context of human contact in prejudice. Hence, the need for religious orientation measure to effectively assess Muslims' contact, relationships and prejudice and be used for further empirical measures.
Notes
1. The idea of extrinsic and intrinsic religious orientation by Allport (1950) was extensively discussed by researchers like Batson, Schoenrade and Ventis (1993), Benson (1988) and Poloutzian (1996). 2. The Qur'an states: "O mankind, We have created you from a male and a female and have made you into nations and tribes for you to know one another. Truly, the noblest of you with God is the most pious" (Qur'an 49:13). 3. Prophet Muhammad is reported to have said: "Allah does not look at your outward appearance and your goods. He looks only at your hearts and your deeds" (Qaysarani 1995, 1:599). 4. It is reported that a companion once asked the Prophet Muhammad: "Who is a neighbour?" He answered: "Your neighbours are forty houses ahead of you and forty houses to your back, and forty houses to your right and forty houses to your left" (Sakhawi 1993, 204). Prophet Muhammad is reported to have said: He who closes his door from his neighbour in fear for his family and wealth is not a believer, and neither is he whose neighbour is unsafe from him a believer. Do you know the right of the neighbour? When in need of help, you support; when in need of a loan, you assist; when impoverished, you provide; when taken sick, you visit; when blessed with goodness, you wish well; when struck with calamity, you condole; when they pass on, you accompany their funeral.
If you have purchased fruit, then make a gift of it to your neighbour. If you have not done so, then take care to conceal it upon your return. Do not allow your child out with it, for it may embitter the neighbour's child or make them covetous. Do you know the right of the neighbour? By the one in whose hand is my soul, they who fulfill the rights of their neighbours do not exceed the few graced by Allah's mercy. He continued to advise them over neighbours, until they thought they would come to include them in their wills. (Bayhaqi 2003, 7:3136) 5. Prophet Muhammad is reported to have said: "The worst food is the food of a wedding banquet in which the rich are invited but the poor are left out" (Ibn al-Mulaqqin 2004, 8:10). The Bible also conveys similar meaning like: "But when you host a banquet, invite the poor, the crippled, the lame and the blind" (Luke 14:13). See Bible Hub retrieved from: http://biblehub.com/luke/14-13.htm (accessed 23 December 2016). 6. Prophet Muhammad is reported to have used the following supplication: "O Lord, I seek refuge with You from leading others astray, causing others to slip or being caused to slip by others, doing wrong or being wronged by others, or behaving foolishly or being treated foolishly by others". Sulayman (1969, 5094).
7. In a famous decree, Abu Bakr al-Siddiq, the first Caliph, told his military commander: Do not commit treachery or deviate from the right path. You must not mutilate dead bodies; do not kill a woman, a child, or an aged man; do not cut down fruitful trees; do not destroy inhabited areas; do not slaughter any of the enemies' sheep, cow or camel except for food; do not burn date palms, nor inundate them; do not embezzle (e.g., no misappropriation of booty or spoils of war) nor be guilty of cowardliness … You are likely to pass by people who have devoted their lives to monastic services; leave them alone. (Ibn Kathir 1995, 2:320) 8. Prophet Muhammad is reported to have said: "No contagion, no pessimism, no haammah [bird] and no Safar" (Albani n.d., 783). The pre-Islamic Arabs used to think that whenever this bird landed on anyone's house, somebody who lived in that house would definitely die. As for Safar is a worm which used to dwell in the bodies of some animals (as a disease) and that this disease was contagious. 9. Prophet Muhammad is reported to have said: "Five are from the natural practices: circumcision, shaving the pubic hair, cutting the moustache short, clipping the nails, and plucking the armpit hairs" (Bukhari 1980, 5889). Prophet Muhammad is reported to have also said: "Keeping your eyes down, clearing the streets of obstacles, responding to Salam greetings, enjoining virtuous deeds and forbidding evil" (Bukhari 1980(Bukhari , 2465. In another narration, he added: "Helping the aggrieved and guiding the aberrant" (Bazzar Ahmad n.d., 1:472). Prophet Muhammad is reported to have said: "Allah is Pleasant and loves pleasant things, Clean and loves cleanliness, Generous and loves generosity" ('Asqalani 2001, 4:254). He is also reported to have said: "Anyone offered rayhan (basil perfume) should not decline it. It is light in weight and fragrant in scent" (Sulayman 2001, 516). 10. Such as the following report on Prophet Muhammad: "Whosoever is slowed down by his deeds will not be hastened forward by his lineage" (Sulayman 1969, 3643). 11. The tradition of Prophet Muhammad illustrates this point well: "Women have certain rights over you and you have certain rights over them" (Ghazali 1997, 454). Prophet Muhammad is also reported to have said: "Treat women nicely" (Muslim 1954(Muslim , 1468. In another narration, Prophet Muhamad is reported to have said: "Women are men's partners" (Sulayman 1969, 236). 12. See Bensaid and Grine (2014, 13). Prophet Muhammad is reported to have said: "God has left no excuse for the person who lives to be sixty or seventy years old; God has left no excuse for him; God has left no excuse for him" (Bayhaqi 1992, 3:370). 13. One may draw on a number of traditions driving much of the Muslim public ethical norms towards the elderly such as: "He is not one of us who does not show mercy to our young ones and esteem to our elderly" (Ibn Hibban 1993, 458). 14. Prophet Muhammad is reported to have said: "No Muslim is pricked with a thorn, or anything larger than that, except that a good deed will be recorded for him and a sin will be erased as a reward for that" (Muslim 1954(Muslim , 2572. 15. Prophet Muhammad is reported to have said: "(Allah) will reward the one whose two dear things (that's his eyes) were taken away from him with Paradise" (Tirmidhi n.d., 2401).
16. Prophet Muhammad is reported to have left Ibn Umm Maktum twice as his successor in Madinah to lead the prayer, though he was blind (Ahmad). Ibn Umm Maktum was a muezzin of Allah's Messenger though he was blind (Tabarani 1994, 1:6). | 9,098.6 | 2018-01-01T00:00:00.000 | [
"Philosophy",
"Sociology"
] |
A High-Throughput Comparative Proteomics of Milk Fat Globule Membrane Reveals Breed and Lactation Stages Specific Variation in Protein Abundance and Functional Differences Between Milk of Saanen Dairy Goat and Holstein Bovine
Large variations in the bioactivities and composition of milk fat globule membrane (MFGM) proteins were observed between Saanen dairy goat and Holstein bovine at various lactation periods. In the present study, 331, 250, 182, and 248 MFGM proteins were characterized in colostrum and mature milk for the two species by Q-Orbitrap HRMS-based proteomics techniques. KEGG pathway analyses displayed that differentially expressed proteins in colostrum involved in galactose metabolism and an adipogenesis pathway, and the differentially expressed proteins in mature milk associated with lipid metabolism and a PPAR signaling pathway. These results indicated that the types and functions of MFGM proteins in goat and bovine milk were different, and goat milk had a better function of fatty acid metabolism and glucose homeostasis, which can enhance our understanding of MFGM proteins in these two species across different lactation periods, and they provide significant information for the study of lipid metabolism and glycometabolism of goat milk.
INTRODUCTION
Milk fat is an important fraction of milk synthesized by the endoplasmic reticulum of mammary epithelial cells. It is a droplet composed of a neutral triglyceride core wrapped by a thin trilayered membrane. Milk fat globules containing cytoplasmic components are retained between the membrane layers as the fat droplets are released into milk. Therefore, milk fat is mainly composed of cholesterol, polar lipids, neutral lipids, and a protein group from the membrane and cytoplasmic crescents (1,2). Protein, which accounts for 22-70% of the MFGM matter, not only provides protection to core milk fat but also has a series of biological functions, such as preventing infection of enteric pathogens, promoting immune and neurological functions, as well as the development of newborns (3,4).
Due to the health benefits and nutritional values, MFGM proteins have attracted growing attention in dairy products. The major proteins in MFGM include lactadherin, mucin one, xanthine dehydrogenase/xanthine oxidase, fatty acid synthase, fatty acid-binding proteins lipophilin, and butyrophilin with physiological functions (5). For example, xanthine dehydrogenase/xanthine oxidase, one of the main MFGM enzyme proteins, has been reported to reveal antimicrobial properties and immuno-protective function. Xanthine dehydrogenase/xanthine oxidase in breast milk reacts with infant saliva to produce an effective combination of irritant and inhibitory metabolites that regulate the gut-microbiota (6). Lactadherin (milk fat globule-EGF factor 8) is a peripheral glycoprotein from human milk, which promotes the clearance and phagocytosis of apoptotic cells, and regulates the immune response (7).
With the accelerated development of proteomics technology in recent years, a large number of proteins have been identified and quantified in the milk of bovine (8), buffalo (9), donkey (10), and other mammals (11). Among them, because bovine milk is the major substitute for human milk, the comparative proteomics between bovine and human milk as well as bovine and human milk were extensively studied. Compared with bovine milk, goat milk-based dairy products may be less allergenic and more easily digested to the infants (12). In view of the unique economic significance, more and more studies have been absorbed in the nutritional and protective properties of proteins in goat milk. Chen et al. have studied the heat-dependent changes of goat milk protein, and found out that heat processing can improve protein digestibility, which was conducive to anti-atherosclerosis therapy (13). They also investigated the protein changes of goat milk during homogenization. The results showed that the homogenized goat milk proteome has changed significantly, which was mainly related to glycolysis/gluconeogenesis metabolism (14). These studies extend the understanding of protein composition in different processes. Major MFGM proteins of goat milk have been reported, which are significantly different from that of bovine. Sun et al. have characterized and compared the MFGM proteins of both Guanzhong goat and Holstein cow milk, using proteomic techniques (15). Furthermore, they analyzed and compared the MFGM proteomes of colostrum and mature milk of Xinong Saanen goat milk (16). However, the MFGM proteome is also affected by species.
Despite the poorly worldwide production of goat milk compared with the bovine, in the past years, there has been more and more interest in the in-depth characterization of its protein composition. This analysis was focused on MFGM proteins from two mammals (bovine and goat) and different lactation periods (colostrum and mature milk) to characterize the composition in conjunction with biological activity, localization, and molecular function of MFGM proteins differences related to lactation. The purpose was to reveal the differences of nutritional value and physiological states of these two species across different lactation periods to provide potential directions for infant formula and functional food development, as well as expand our current knowledge of MFGM proteome.
Sample Collection
The sample collection and preparation were shown in Supplementary Figure 1. The samples were collected, followed by the method reported by Sun et al. (16). The samples were collected at the Holstein bovine and Saanen dairy goat farm in Xi'an, Shaanxi province, China. Ten bovine colostrum (0-5 days postpartum), 10 mature-milk (1-6 months postpartum), 10 goat colostrum (0-5 days postpartum), and 10 mature-milk (1-6 months postpartum) samples were obtained from 20 healthy bovines and 20 healthy goats in the first lactation. All of the 40 animals were aged between 1 and 4 years old, and the animals of each species were under identical environmental conditions. Each sample of bovine and goat milk was collected twice a day and then mixed to dispel the effect of the sampling time of milk samples. These samples were transported to the laboratory on ice and stored at −80 • C. Ten milk samples of each group were mixed to refrain from the influence of individual differences on MFGM protein in various lactation stages before the analysis. All handling practices involving animals carefully followed all the recommendations of the Directive 2010/63/EU of the European Parliament for the protection of animals for scientific purposes.
The extraction of MFGM proteins was conducted as described by Lu et al. (17) with minor modifications. Briefly, 50-ml milk samples were centrifuged at 12,000 × g for 40 min at 4 • C. The supernatant (top layer) was transferred to another centrifuge tube and washed three times at 25 • C for 10 min, with 0.1 mol L −1 PBS (pH 6.8) and centrifuged at 10,000 × g for 15 min at 4 • C subsequently to remove residual whey proteins and caseins. Then we washed the cream twice, using ultrapure water to dislodge the residual salt ions. Finally, 0.4% SDS (1:1, v/v) was added to dilute the washed cream, sonicated for 1 min and centrifuged at 10,000 × g at 4 • C for 40 min to separate the fat fraction. The MFGM proteins were collected in the aqueous phase (bottom layer), and their concentration was measured by BCA assay (Thermo Scientific Pierce BCA protein assay kit, USA).
Protein Digestion
The MFGM protein was reduced, alkylated, and digested, and followed the method reported by Lu et al. (18). For each milk sample, three independent biological replicates were made. First, 10-µL MFGM protein was dissolved in 100-µL, 50-mmol L −1 NH 4 HCO 3 , and then 10-µL, 100-mmol L −1 dithiothreitol was added and incubated at 56 • C for 30 min. Subsequently, the MFGM protein was alkylated with 15 µL of 55-mmol L −1 iodoacetamide in dark for 30 min at room temperature and then adding sequencing grade-modified trypsin to digest the MFGM protein at a ratio of 1:100 enzyme/protein for 16-18 h at 37 • C and terminated the reaction by adding 1% formic acid. Finally, the peptides mixture was desalting by Oasis HLB cartridges (Waters Cooperation, Milford, MA, USA), dried by a vacuum centrifuge and then resuspended in 40 µL of 0.1% (v/v) formic acid.
Liquid Chromatography Tandem Q-Orbitrap Mass Spectrometry
Peptide separation was performed by EASY-nLC 1000 system (Thermo Scientific, San Jose, CA), equipped with a C18-reverse phase column (75-µm inner diameter, 10-cm long, 3-µm resin; Thermo Scientific) at 200 nL/min and 35 • C for a total run time of 100 min. Solution A (0.1% formic acid in water) and solution B (0.1% formic acid in 80% ACN) were used as eluents for the peptide separation according to the following elution gradients: 5-35% solution B for 50 min; 35-100% solution B for 25 min, followed by 15-min washing with 100% solution B, and return to 5% B in 0.1 min, re-equilibration during 9.9 min with 5% solution B.
The peptide eluted from the column was ionized by a Q-Exactive (Thermo Fisher Scientific, Waltham, USA) mass spectrometer in a positive mode. The spray voltage was operated at 3.8 kV. The m/z scan range of single MS scans of peptide precursors was m/z 300-1,700 at a resolution of 70,000 (at m/z 200). The AGC (automatic gain control) was 3 e 6 , and the maximum injection time was 200 ms. The top 20 most intense precursor ions with charge ≥2 determined by MS scan were used to obtain MS/MS data at a resolution of 17,500 by using a higher normalized collision energy of 27 eV. The AGC was 1 e 5 and the maximum injection time was 50 ms. In order to avoid superfluous fragmentation, the dynamic exclusion time was set to 30 s.
Data Analysis
The raw LC-MS/MS files with three replicates were obtained for MFGM proteins of each milk group. Two proteins identified from three biological replicates of each milk group were used for subsequent analysis. The data analysis was carried out using the MaxQuant software (Max Planck Institute of Biochemistry, Martinsried, Germany, version 1.6.7.0), with Andromeda as a peptide search engine (Matrix Science, version 2.4), and searched against the database of Caprinae (67,040 entries, 02/08/2019) and Bos taurus (64,796 entries, 02/08/2019) organism group with reverse sequences generated by MaxQuant.
Search parameters were a first search peptide mass tolerance of 20 ppm and main search peptide mass tolerance of 4.5 ppm. Methionine oxidation and protein N-terminal acetylation were defined as variable modification and carbamidomethyl of cysteine defined as a fixed modification for both identification and quantification. A trypsin/P was adjusted as a proteolytic enzyme with a maximum of two missed cleavages. A maximum of 0.01 false discovery rates (FDRs) and at least two unique peptides for each protein were demanded for reliable identification and quantification. Label-free quantification (LFQ) was enabled in MaxQuant.
The MFGM proteins identified in at least two of the three biological replicates were used for subsequent statistical analysis. The LFQ intensity of identified proteins was analyzed using one-way ANOVA test. MFGM proteins with p < 0.05 and fold change >2 in the relative abundance ratios were considered to be differentially expressed MFGM proteins (DEMPs) between milk groups. Principal component analysis (PCA) and partial least squares regression-discriminant analysis (PLS-DA) were used to construct the recognition model and prediction models, respectively, among milk groups by using SIMCA 15 (19).
Bioinformatics Analyses
The molecular function, cellular component, and biological process of all identified MFGM proteins according to their gene ontology (GO) annotations and Kyoto Encyclopedia of Genes and Genomes (KEGG) pathway analysis were achieved, using DAVID Bioinformatics Resources 6.8 (https://david.ncifcrf. gov/). Conversion of the genes of MFGM proteins was using Retrieve/ID mapping (https://www.uniprot.org/uploadlists/). Protein-protein interaction (PPI) network construction was performed, using STRING (https://string-db.org/), with the DEMPs displayed by proteomic data used as input. Statistical differences were proclaimed significantly if p ≤ 0.05.
Component Analysis of MFGM Proteins From Different Milk Groups
In this study, 331 MFGM proteins in goat colostrum (GC), 250 in goat mature milk (GM), 182 in bovine colostrum (BC), and 248 in bovine mature milk (BM) were identified and quantified using LC-Q-Orbitrap mass spectrometer (Supplementary Table 1). These proteins spanned more than six orders of magnitude. As shown in Figure 1, there were 136, 72, 37, and 78 uniquely expressed MFGM proteins identified in goat colostrum, goat mature milk, bovine colostrum, and bovine mature milk, respectively. The uniquely expressed proteins in goat colostrum included calreticulin, cofilin-1, and methanethiol oxidase in goat mature milk included FA complementation group I, phosphatidylethanol-amine binding protein 4, and perilipin-1. Approximately, 45 MFGM proteins, coupled with cream fractions, were identified in four milk groups, including the complement component 3, fatty acid synthase, sodiumdependent phosphate transport protein 2B, and pericentrin, indicating that goat milk was a substitute for bovine milk to some extent in the research and development of infant milk powder and functional food, and also providing an orientation for the further functional development of goat milk.
In addition, proteins from caseins and whey, such as αs1-casein, αs2-casein, β-casein, κ-casein, lactoferrin, and β-lactoglobulin, were also identified in the MFGM fractions, which could be owing to the residual contamination of the proteins during MFGM extraction (20). In our study, polymeric immunoglobulin receptor (PIGR) was highly abundant in four milk groups. PIGR binds with polymerized immunoglobulin A and immunoglobulin M and transport them to perform immune functions across cell membranes (21). The other identified immunoglobulin proteins included IGL@ protein (IGL@), IGK protein (IGK), immunoglobulin heavy-constant mu (IGHM), and immunoglobulin J chain (JCHAIN). Previously observed major proteins, such as lactadherin, butyrophilin, and xanthine dehydrogenase/oxidase, were conserved in MFGM across four groups, which indicate the robustness of the methodology.
Chemometrics of MFGM Protein in Various Lactation Periods
Before chemometric analysis of MFGM protein, preprocessing of LC-MS/MS data was implemented. The correct normalization of each milk group can eliminate the systematically differences between features. However, the total peptide ion signals necessary to perform LC-MS/MS were distributed over several adjacent runs. Therefore, it was necessary to know the normalization coefficient (N) of each fraction to sum the peptide ion signal. Based on the least overall proteome variation, the quantities of proteins can be determined via a global optimization procedure after the intensities were normalized to a normalization factor as free variables. Hence, in sample A, the total intensity of a peptide ion P was defined as where the index k covered all isotope patterns of peptide ion P in sample A.
A triangular matrix containing all paired protein ratios between any two samples was constructed. This matrix corresponds to the overdetermined system of equations for the protein abundance distributions in the sample. A subsequent least-squares analysis was performed to reconstruct the abundance profile based on the sum of squared differences in the matrix via the optimal satisfaction of individual protein ratios.
Then the whole profile was rescaled to the cumulative intensity of the samples, thereby retaining the total summed intensity of the protein over all samples, which was the "LFQ intensity" (22).
Chemometric analysis is a technique involving statistical methods to comprehend chemical information generated by analytical instruments. To evaluate whether the data of proteomic analysis can be engaged to visually differentiate the four group milk samples, PCA analysis as an unsupervised data analysis was carried out by loading the LFQ intensities of all detected proteins as variables and the different MFGM matrices as observation points. In this case, the data from all the replicates were used. A clear separation of all MFGM matrices can be observed on the PCA score plot (Figure 2A) based on the first two principal components. Samples close to each other on the PCA score plot revealed similar properties, while the milk samples far from each other revealed dissimilar molecular weight in protein mass spectrographic analysis. The first two principal components accounted for 81.2% [R 2 X (1) + R 2 X (2)] of the total variation in the data. No outlier was observed by ellipse Hotelling's T2. The first principal component (45% of the total variation) clearly divided the samples into four wellseparated clusters, corresponding to the four groups, showed that the distribution differences among the goat colostrum, goat mature milk, bovine colostrum, and bovine mature milk were mainly due to biological reasons. Meanwhile, R 2 X (cum) was cumulative R 2 X up to the specified component, where R 2 X was the fraction of X variation modeled in the component. Q 2 (cum) was the cumulative Q 2 up to the specified component, where Q 2 was overall cross-validated R 2 X. R 2 (cum) and Q 2 (cum) were the critical parameters to evaluate the quality of the PCA model, which, respectively, reflect the degree of interpretation of the principal components to X variables and the predictive ability of the model. Figure 2A showed that R 2 X (cum) was 0.812 and Q 2 (cum) was 0.966, illustrating that the PCA model was stable and predictable.
PLS-DA is a chemical projection method that associates X and Y variable blocks via a linear multivariable model. The objective is to find the direction in X space that divides the classes according to the sample set with known class members (23,24). In the built PLS-DA model, R 2 X (cum) was 0.812, R 2 Y (cum) was 1, and Q 2 (cum) was 1, which denoted that the PLS-DA model has good fitting ability and prediction performance ( Figure 2B). VIP scores were shown in Figure 2C, and variables with VIP scores higher than 1 were considered to have significantly contributed to the model. To evaluate the stability and reliability of PLS-DA model, cross-validation was adopted. The total correct classification rate of four MFGM proteins was 100%. The statistical significance of the predictive quality parameters in the built PLS-DA model was validated by 200 permutation tests (Figure 3). The Y-intercept of R 2 and Q 2 was 0.165 and −0.711, respectively, which ensured that the PLS-DA model was not overfitting.
DEMPs, Respectively, in Colostrum and Mature Milk of Goat and Bovine
Three hundred and thirty-one MFGM proteins were identified in goat colostrum and 182 in bovine colostrum. T-test and fold change analysis were used to analyze the difference of the MFGM proteins in colostrum of goat and bovine, and the MFGM protein with p < 0.05 and at least two-fold was considered as the cutoff criteria of differential expression. As shown in Table 1, among 74 common proteins in colostrum of goat and bovine, 49 were differentially expressed. We found out that there were 22 upregulated and 27 downregulated MFGM proteins in goat colostrum. The levels of sodium/nucleoside cotransporter, lipoprotein lipase, sodium-dependent phosphate transport protein 2B, xanthine dehydrogenase/oxidase, and fatty acid synthase were higher in goat colostrum, while the levels of cathelicidin-1, lipopolysaccharide-binding protein, alpha-enolase, and vitamin D-binding protein were higher in bovine colostrum. Two hundred and fifty MFGM proteins were identified in goat mature milk and 248 in bovine mature milk. As shown in Table 2, among 90 common proteins in mature milk of goat and bovine, 63 were differentially expressed on the basis of p < 0.05 and at least two-fold, with 32 upregulated and 31 downregulated MFGM proteins. The levels of lipoprotein lipase, fatty acid synthase, sodium-dependent phosphate transport protein 2B, and apolipoprotein E were higher in goat mature milk, while the levels of lipocalin 2, sodium/nucleoside cotransporter, apolipoprotein A-I, and polymeric immunoglobulin receptor were lower in bovine mature milk than in goat. The abundance of lipoprotein lipase and fatty acid synthase was higher in goat mature milk, which was the same as that in colostrum. Lipopolysaccharidebinding protein was higher in bovine mature milk in contrary. Lipocalin-2 can bind and eliminate enterochelin, a high-affinity siderophore, to reduce the accessibility of bacteria to iron and inhibit its growth, and participate in the modeling of immune response. Apolipoprotein A-I is a kind of lipoprotein expressed by glial cells, which is strongly induced in aging, injury or neurodegeneration, involved in the peripheral metabolic regulation and lipid processing of chylomicron, the occurrence and development of Parkinson's disease and Alzheimer's disease, and may play a neuroprotective role in the brain (16).
GO, KEGG Pathway, and PPI Analysis of the DEMPs
For the purpose of comparing the biological functions of MFGM proteins in colostrum, 49 DEMPs in colostrum of goat and bovine were analyzed by gene ontology (GO) functional annotation, which was divided into three categories of molecular function (MF), cellular components (CC), and biological process (BP). This was useful to further understand the colostrum MFGM proteins functions in goat and bovine. The most significant enrichment annotation information (p < 0.05) in each branch was shown in Figure 4A, in which the prevalent biological processes were a response to dehydroepiandrosterone, response to 11-deoxycorticosterone, response to estradiol, and response to progesterone. Others were involved in negative regulation of endopeptidase activity, acute-phase response, lactose biosynthetic process, and cell adhesion. The MFGM proteins were highly enriched in extracellular space, extracellular exosome, and blood microparticle, which illustrated that the majority of these proteins present due to leakage of the protein from the blood serum into the milk at the tight junctions in the cells in the mammary gland. Other enriched origin categories were Golgi lumen, Golgi apparatus, and extracellular region. Three prominent molecular functions of DEMPs in colostrum of goat and bovine were protein binding, transporter activity, and cytoskeletal protein binding, which was consistent with the result of Cunsolo et al. (25). Lactose synthase activity and structural molecule activity were also significantly represented.
The most significant enrichment annotation information (p < 0.05) in each branch of mature milk was shown in Figure 4B, in which the prevalent biological processes were response to dehydroepiandrosterone, response to 11deoxycorticosterone, response to progesterone, and response to estradiol. It was consistent with the biological processes involved in colostrum, indicating that these biological processes occurred in the whole lactation period. In stressful situations, the adrenal cortex reacted to ACTH and began to secrete dehydroepiandrosterone that had been proved to exert antiinflammatory and antioxidant effects, and played a protective and regenerative role (26). Others were involved in very-low-density lipoprotein particle remodeling, triglyceride catabolic process, and cholesterol biosynthetic process. The MFGM proteins were highly enriched in extracellular exosome, blood microparticle, and extracellular space, which was similar to the DEMPs in colostrum of goat and bovine. Other enriched origin categories were Golgi lumen, membrane, and chylomicron. The prominent molecular functions were structural molecule activity, heparin binding, and protein binding.
A KEGG pathway was employed to analyze the main pathways of MFGM protein differentially expressed in colostrum of goat and bovine. As shown in Figure 5A, the DEMPs mainly involved in the pathways of the phagosome, a PPAR signaling pathway, proteoglycans in cancer, ECM-receptor interaction, and galactose metabolism. As shown in Figure 5B, the DEMPs in mature milk mainly involved in the pathways of viral myocarditis, a PPAR signaling pathway, salmonella infection, and fatty acid biosynthesis.
PPI analysis of DEMPs in the colostrum and mature milk of goat and bovine was implemented to obtain a colorcoded network, revealing the correlation between DEMPs (Supplementary Figure 2). The final network of colostrum consists of 49 nodes (proteins) and 99 edges (interactions). The avg. local clustering coefficient was 0.533, and the PPI enrichment p-value was lower than 1.0 e −16 , indicating that the 49 DEMPs were biologically connected. Most highly interacting protein nodes in colostrum were divided into five communities, including a cellular process, cell adhesion, immune response, lactose synthase activity, and a metabolic process. Similarly, the final network of mature milk consists of 63 nodes (proteins) and 148 edges (interactions). The avg. local clustering coefficient and p-value were 0.575 and lower than 1.0 e −16 , respectively, indicating that the 63 DEMPs were biologically connected. Furthermore, these highly interacting protein nodes were divided into three communities, including a cellular process, cell adhesion, and response to extracellular stimulus. Consequently, in light of the bioinformatics analysis of proteomics, the MFGM proteins in colostrum might be led to the intervention of related health issues, while the bioactivities of MFGM proteins between goat and bovine might be different.
Protein Profile Analysis Based on Chemometrics
A label-free proteomic approach was employed to identify the MFGM proteins in goat and bovine milk. Compared with the previous reports, 102 MFGM proteins, mainly calciumbinding protein, cytoskeletal protein, and intercellular signal molecule were added in bovine milk, 114 in goat milk, mainly a metabolite interconversion enzyme, translational protein, cytoskeletal protein, and nucleic acid-binding protein [Supplementary Table 2; (15)(16)(17)27)]. The nucleotide-binding proteins in this study were mainly DNA helicase, followed by several members of Ras super family (Rhos and Rabs). All these proteins are likely involved in the secretion of secreted milk components by vesicle (17). Furthermore, the results of PCA, PLS-DA, and VIP analysis found out that nucleobindin-1 (NUCB1), folate receptor alpha (FOLR1), vitamin D-binding protein (GC), thrombospondin-1 (THBS1), and beta-1,4galactosyltransferase 1 (B4GALT1), with the lower abundance in goat mature milk than in other three milk groups, can be used as markers to distinguish the four milk groups. NUCB1 is Golgi-localized soluble protein, which contains multiple putative functional domains. NUCB1 localized in extracellular was considered to be a regulator of matrix maturation in bone (28,29). FOLR1 can internalize folates into the cells, which is crucial to DNA repair and synthesis, and mediate the activation of pro-oncogene STAT3, which contributes to angiogenesis, tumor proliferation, and metastasis (30). THBS1 has been demonstrated to participate in mechano-signal transduction and is specific at the level of apoptosis induction (31). GC has the metabolic effects of influencing bone metabolism, chemotaxis, actin scavenging, innate immunity, modulation of inflammatory processes, and binding of fatty acids (32). As mentioned above, the data from our study not only provide a comprehensive understanding of the MFGM protein compositions among the four milk samples but also reveal the differences of MFGM proteins among different species of mammals. The results exhibited a scientific basis for the development of functional products, using goat milk.
A number of MFGM proteins were obviously differentially expressed in the goat and bovine milk, whose functions may associate with the flavor and protection of the lamb or calf from infections. Milk fatty acids are extracted from the arterial blood or synthesized de novo in the mammary gland and involve many kinds of mammary enzymes, including lipoprotein lipase and fatty acid synthetase (33). The abundance of lipoprotein lipase was higher in colostrum of goat milk than that of bovine milk. Lipoprotein lipase was more closely bound with fat globules and has a better correlation with spontaneous lipolysis in goat milk. It can release fatty acids from lipoproteins and chyle particles, which may relate to the differences in the flavor of goat and bovine milk. Zhu et al. (34) have clarified that inhibition at the gene-expression level of fatty acid synthetase restrained the accumulation of TAG and the formation of lipid droplet by reducing esterification and lipogenesis and promoting lipolysis in goat milk. In addition, fatty acid synthetase also helped generate toll-like receptor 4 (TLR4), which is vital for lipid metabolism regulation. Lipopolysaccharide-binding protein, a 58-60 kDa protein, catalyzed the transfer of bacterial lipopolysaccharide to CD14, which exists in soluble form and facilitates lipopolysaccharide presentation to TLR4 as a cell surface receptor. This activates the intracellular signaling pathways and promotes the upregulation of adhesion molecules and proinflammatory cytokines, which are participated in the innate immune response (35). The low abundance of lipopolysaccharide-binding protein in goat colostrum may be related to its low sensitization. Cathelicidin-1 found in epithelial and neutrophils cells is an antimicrobial peptide and has profound impacts on wound healing, inflammation, and the regulation of adaptive immunity (36). The existence of antimicrobial proteins in the colostrum of goat and bovine milk explained that the protection from milk was necessary for newborn mammals and varies across species.
DEMPs Analysis Based on Bioinformatics
Bioinformatics analysis showed that the response to dehydroepiandrosterone, 11-deoxycorticosterone, estradiol, and progesterone were the most abundant biological processes in colostrum and mature milk. Previous researches showed that the interaction between low-density lipoprotein and heparin results in irreversible structural changes in apolipoprotein B and affects low-density lipoprotein oxidation, phospholipolysis, and fusion (37). Dehydroepiandrosterone, 11-deoxycorticosterone, estradiol, and progesterone are endogenous steroid hormones, which are often conjugated with proteins, secreted by endocrine glands and endocrine cells, dispersed in other organs, and then transported to various organs through blood, provided with the function of coordinating and controlling tissue, organ metabolism, and physiological function (38). Progesterone is a good marker to determine the milk production function status, which may affect the milk production of goats and make it lower than that of bovines. Cholesterol-rich plaques accumulated in the arteries, preventing enough blood flow to the heart, causing cardiovascular disease, bring about many medical suggestions, and require a reduction of the cholesterol intake (39). However, throughout the lactation periods, the cholesterol levels in milk decreased sharply, and the cholesterol concentration in mature milk of goat (11.64 ± 1.09 mg/dL) was lower than that in bovine (20.58 ± 4.21 mg/dL), which is consistent with the results of the GO analysis of the DEMPs in mature milk of goat and bovine. The result of the most abundant biological processes in colostrum and mature milk indicated that goat milk is more comfortable for the human body. The extent of a triglyceride catabolic process, which occurs via acid lipolysis in the lysosome and neutral lipolysis in the cytoplasm, can regulate lipid droplet size (40). LPL and APOE as the critical enzymes that participate in the triglyceride catabolic process were higher in mature milk of goat than that in bovine goat milk, which confirmed that the regulating lipid metabolism in MFGM proteins of the goat was superior to that of bovine MFGM proteins.
The phagosome pathway is the important process of catabolism and transport, which has been reported in the Guanzhong goat and Holstein cow mature milk (16). Yang et al. (41) revealed that the phagosome pathway was a complicated process of organic uptake and elimination of apoptotic cells and pathogens, which contributes to inflammation, host defense, and homeostasis in tissues. Thus, the MFGM proteins may be crucial to the immune system of newborn mammals. The ECM-receptor interaction pathway has been found to be associated with depotspecific adipogenesis in cow and overrepresented in specific cattle breeds related to the adaptive immune response after virus inoculation in Holstein cattle (42). The structural functional diversity of proteoglycans makes them as key mediators of the interaction between tumor cells and host microenvironment, and directly participates in the tissue and dynamic remodeling of the milieu. As constituents of the ECM or extracellular milieu, proteoglycans may invariably participate in the control of a variety of oncogenic events in a multivalent manner (43). KEGG pathway analysis also displayed that the differential MFGM proteins in colostrum goat and bovine significantly regulated glycometabolism through the pathway of galactose metabolism and revealed that the MFGM proteins of goat and bovine milk possessed different effects on glucose metabolism. Peroxisome proliferator-activated receptors (PPARs) are the members of the nuclear hormone receptor superfamily, which has three member isotypes: PPARα, PPARβ/δ, and PPARγ, and is ligandactivated transcription factors. PPARs govern the expression of the crux molecules in fatty acid metabolic pathway, including the absorption, oxidation, and storage of fatty acids. Meanwhile, PPARγ maintains glucose homeostasis by activating glucose transporter 2 and glucokinase in the pancreas and liver (44). LPL, ACSL1, FABP3, and CD36 play key roles in the PPAR signaling pathway, and the abundance of these proteins in goat milk is higher than that in bovine milk, which indicated that goat milk has better function of fatty acid metabolism and glucose homeostasis. Due to the different functions of MFGM protein, the MFGM protein expressed in different lactation periods of bovine and goat provided significant information for functional food and infant formula. According to the nutritional needs of different lactating infants, functional proteins can be added to the corresponding formula, which is conducive to the health of infants.
In a previous study, the efficiency of preventing a series of enteric inflammatory and infectious diseases by milks has been documented. Viral myocarditis and salmonella infection are the subcategory of cardiovascular disease and infectious disease, respectively. The MFGM protein in mature milk significantly inhibited the internalization and binding of salmonella during the growth of mammals, and the inhibitory effect of goat milk was stronger than that of bovine milk, which was in line with our data (45). Sun et al. (15) reported that most of goat MFGM proteins were related to metabolism processes, including lipid metabolism and carbohydrate metabolism. According to the KEGG analysis results, the difference of fatty acid synthesis pathway between bovine and goat could mainly be a consequence of fatty acid synthetase activity differences between bovine and goat. Fatty acid synthases are located on the surface of endoplasmic reticulum and on the cytoplasmic lipid droplets in mammary epithelial cells. The endoplasmic reticulum possesses a series of membrane enzymes, such as palmitoyl-CoA and glycerol 3-phosphate, which can synthesize triglycerides from activated fatty acids. Palmitoyl-CoA, a major product of fatty acid synthases, and fatty acid synthases were directly connected in the fatty acid synthesis pathway (46). In addition, glucose and fructose promote the accumulation of triglycerides in milk to convert sugar to fat, the allosteric inhibition of fatty acid oxidation via increased availability of triose phosphate precursors and acetyl-CoA and metabolites, such as malonyl-CoA, for fatty acid formation via glycerol-3-phosphate biosynthesis and de novo lipogenesis (47). Regulation of fatty acid biosynthetic could inhibit hepatic steatosis and lipid accumulation.
CONCLUSIONS
Our study described a more specific strategy to provide insights into proteome differences between GC, GM, BC, and BM. Bioinformatics analysis displayed that these DEMPs in colostrum significantly regulated glycometabolism through the pathway of galactose metabolism, and the DEMPs in mature milk regulated lipid metabolism through the pathway of fatty acid biosynthesis. These trials and results could reveal the differences of nutritional value and physiological states between bovine and goat at various lactation periods, and provide direction for the application of goat milk in infant formula food and functional food. Further study on the exact role of these significant difference proteins and the specific mechanism of activating the regulation pathway of galactose and lipid is necessary and could potentially help to determine the new biomarkers or establish the optimized formula of goat milk-based functional food.
DATA AVAILABILITY STATEMENT
The original contributions generated for the study are publicly available. This data can be found here: https://zenodo.org/record/ 4638812#.YHbYPOhKiUm.
AUTHOR CONTRIBUTIONS
WJ: conceptualization, methodology, software, writing-original draft preparation, and supervision. RZ: data curation and writing-original draft preparation. ZZ: reviewing and editing. LS: software. All authors contributed to the article and approved the submitted version. | 7,947.8 | 2021-05-28T00:00:00.000 | [
"Agricultural and Food Sciences",
"Biology"
] |
Simulation of Time-Dependent Schrödinger Equation in the Position and Momentum Domains
The paper outlines the development of a new, spectral method of simulating the Schrödinger equation in the momentum domain. The kinetic energy operator is approximated in the momentum domain by exploiting the derivative property of the Fourier transform. These results are com-pared to a position-domain simulation generated by a fourth-order, centered-space, finite-differ-ence formula. The time derivative is approximated by a four-step predictor-corrector in both domains. Free-particle and square-well simulations of the time-dependent Schrödinger equation are run in both domains to demonstrate agreement between the new, spectral methods and established, finite-difference methods. The spectral methods are shown to be accurate and precise.
Introduction
This paper outlines the development of simulations of the time-dependent Schrödinger equation produced in both position and momentum domains. In the position domain this is the x-y plane. In the momentum domain this is the k x -k y plane, as it is the Fourier transform of the position domain [1]. The simulations demonstrate the accuracy of the spectral methods used in the momentum domain.
All simulations are advanced in time using a four-step predictor-corrector method. The predictor-corrector can be applied independently in both position and momentum domains to step the simulation forward in time. The predictor-corrector is generated using Lagrange polynomials, outlined by [2] and [3]. The predictor formula found here is shown to be consistent with established Adams-Bashforth formulas [3].
The position-domain approximation of the kinetic energy operator is derived using Lagrange polynomials and consistent with results from [4]. In the position domain the approximation to the kinetic energy operator is fourth-order accurate. In the momentum domain, the kinetic energy operator approximation is global-order accurate because it relies on the derivative property of the Fourier transform [5]. The software written to generate these simulations uses the Fastest Fourier Transform in the West (FFTW) to transform between position and momentum domains [6]. Simulating the time-dependent Schrödinger equation in the momentum domain achieves higher orders of spatial accuracy. The performance and precision of momentum-domain simulations is comparable to position-domain simulations.
Given an initial state 0 y at 0 t = , the four-step predictor-corrector requires the creation of n j n y y j t − = − ∆ for 1, , 4 j = in order to compute the first predictor-corrector time-step. A simple backwards Euler method, outlined in [4], [7]- [12], is used to generate the wave function at these early time-steps. Each of these early states for 1, , 4 j = is re-normalized after their creation to ensure minimum initial error. The first simulation is a free particle with no imposed boundary conditions, when the Hamiltonian consists only of the kinetic energy operator. This simulation demonstrates the difference in boundary conditions of each domain. In the position domain, this is equivalent to an infinite square-well potential, or particle-in-a-box. When the wave function reaches one boundary, it is reflected back. In the momentum domain, this is equivalent to periodic boundary conditions. When the wave function disappears into one boundary it will reappear in the opposite boundary, travelling in the same direction. This is to establish a relative performance benchmark when only the kinetic energy operator is applied.
Second, a finite square-well potential of 100 eV is imposed in both domains. This simulation demonstrates the computational burden associated with imposing the same initial and boundary conditions in both domains. Application of the potential function in the position domain is carried out by entry-wise multiplication of the wave function and potential function lattices. In the momentum domain, this operation is equivalent to convolution. Rather than carry out this time-consuming operation, the wave function in the momentum domain is transformed back to the position domain at every time-step in the simulation in order to apply the potential function. The kinetic energy operator is applied when the wave function has been transformed forward into the momentum domain.
Each of the simulations begins with an initial state of a two-dimensional wave packet with a Gaussian envelope. The simulations are stepped forward in time and the complex-valued wave function components and densities, as well as some expectation values, are captured incrementally. The wave function components and densities are converted to image format and animated [13].
Methods
The following subsections outline the numerical methods used to generate solutions to the Schrödinger equation. (1)
Time Derivative
The same four-step predictor-corrector method is used to step the position and momentum domain simulations forward in time. The predictor and corrector start with the basic form of the differential equation
Predictor
The predictor uses the integral form of the differential equation.
The function f is approximated by Lagrange polynomials [2]. Once the polynomial approximation to f has been substituted, the integral is straightforward to compute. This yields Adams-Bashforth coefficients which calculate the predicted value n y [3].
Corrector
The corrector uses the original form of the differential equation.
( )
n n y f y y f t For the corrector, y is approximated by Lagrange polynomials [2]. Once the polynomial approximation to y has been substituted, the first derivative is straightforward to compute. This yields the following coefficients which calculate the value n f .
The predicted value n y is substituted for n y in the function f to get n f .
Application to the Schrödinger Equation
In the position domain, the Schrödinger equation is written as follows using operator notation as shorthand. The Hamiltonian operator Ĥ is written with a P superscript to denote the position domain and with superscript M to denote the momentum domain.
At every time step, the predictor calculation is carried out.
That result is plugged in to the corrector calculation to advance the simulation forward one time-step.
The following sections outline the development of the Hamiltonian operator ˆP H in the position domain and ˆM H in the momentum domain.
Free Particle, Kinetic Energy Operator
The first simulations in the position and momentum domains assume free particle conditions. The particle is given the mass of an electron. Only the kinetic energy operator applies in the Hamiltonian.
Before the simulations are started, the position-and momentum-domain lattices must be constructed. This requires fixing the position-domain step sizes x ∆ and y ∆ as well as the number of columns x N and number of rows y N . It is also helpful to define an origin point, 2 The reversal of the y-direction accounts for the fact that computer storage increments the row index as the row moves down.
The momentum domain lattices are constructed according to the relationship 1 . The high frequencies have been shifted into the negative frequencies. Use of the FFTW library requires applying a phase-shift to the position domain before transforming into the momentum domain if negative frequencies are used instead of high frequencies [6].
Position Domain
The approximation to the kinetic energy operator in the position domain was generated using Lagrange polynomials. This was accomplished by approximating the second derivative in one dimension, as the same formula can be applied to all dimensions. This is a centered-space formula accurate to fourth order, requiring five sample points. In terms of the generalized coordinate q , the sample points and a fixed q ∆ describe the set of sample points centered around 2 q . For a function ( ) , the polynomial approximation to f is substituted and the second derivative is calculated. This procedure yields the following approximation to the Laplacian operator: 16 30 The real and imaginary parts of the wave function R iI Ψ = + are calculated independently, yielding a pair of coupled equations. 2 Denote spatial sample points with subscripts: ( ) I I I I I m x
Momentum Domain
The approximation to the kinetic energy operator in the momentum domain was generated using the transform of the derivative operator. Because the momentum domain is the Fourier transform of the position domain, the derivative operator is transformed as well. In terms of the generalized position-domain coordinate q , let the function ( ) F s be the Fourier transform of the function ( ) and s is the generalized momentum-domain coordinate.
The initial position-domain state Ψ is transformed forward into the momentum-domain state Φ , where ( )( )
Square Well, Potential Energy Operator
A square-well potential was chosen to test the effectiveness and relative performance of the simulations of both domains, when the particle interacts with an electrostatic potential. The particle is again given the mass of an electron. For the purposes of this demonstration, the chosen potential must be high enough to reflect most of the particle off the potential step, back into the region where the potential is zero.
The lattice describing the potential must have the number of columns x N and number of rows y N . Before the simulations are started, it is helpful to choose a well boundary index constant step j and potential step size
Position Domain
The application of the potential operator is straightforward in the position domain. The lattices representing the real and imaginary parts of the wave-function Ψ are multiplied entry-wise by the lattice representing the potential function. The simulation can be stepped forward in time once the approximation to the Hamiltonian has been substituted into the predictor corrector formula.
Momentum Domain
In the momentum domain, application of the potential operator transforms to the convolution operation, and * represents convolution [5]. Rather than carry out this computationally expensive operation, the operator is applied in the position domain. This requires transforming back and forth between the position and momentum domains at every time-step. Use of the predictor-corrector complicates this procedure somewhat, because the potential operator must be applied at the predictor step and the corrector step.
For simplicity and readability the procedure below does not write the state functions Ψ and Φ decomposed into their real and imaginary components.
First, the potential operator is applied in the position domain and transformed to the momentum domain.
The kinetic energy operator is applied to Φ and added to * Φ to produce the predictor Hamiltonian ˆM pred H .
( )( )
Following the example of Equation (9) in Section 2.1.3, the corrector formula is applied to produce the corrected value Φ .
Results
All simulations are run under the same constraints and initial conditions to illustrate similarities and differences in the particle's position over time. The free particle simulations show differences in position due to the differences in boundary conditions between the two domains, despite having the same initial conditions. The square well simulations will show strong agreement in position due to the boundary conditions and initial conditions being the same. Regarding precision, all simulations were shown to be accurate to 8 decimal places. This value is stable for the duration of the simulations and was measured by finding the difference between the current state's density function and 1.
The physical constants of electron mass m , electron charge q and Planck's constant are given in standard units. This helps scale the simulations to real-world dimensions, although the simulated particles are much larger than real electrons to show detail. For the square-well simulations, the electrostatic potential is set to +100.0 eV, although this is also expressed in standard units. The well boundary is established at 20 index units inside the lattice boundaries.
The row and column sizes are set to 256 × seconds. Each simulation is run for 100,000 time steps. The lifetime of each simulation is 0.1 picoseconds. The components, densities and expectation values are measured and recorded every 500 time steps. This increment is referred to as a "frame" in the graphs below and each frame is equal to 16 5.0 10 − × seconds. All initial condition parameters are given in index units and the actual parameters are multiplied by the appropriate lattice step size. The initial conditions set the particle on the +x-axis at index +70.0, relative to the origin point and close to the vertical boundary at the end of the +x-axis. The particle is given positive momentum, which will propel the particle into the boundary. The particle's spatial wavelength λ is set to 10.0 index units. The particle's predicted velocity can be found by the conservation of momentum.
The initial velocity of the particle in all simulations is 5 3.64 10 × meters per second in the +x-direction. The initial momentum x k is therefore 9 π 10 × kilogram-meters per second. In units normalized by 2π , as shown in the momentum domain simulations, the initial momentum is
Free Particle
Based on the initial conditions and predicted velocity, the particle is expected to reach the boundary at time 14 3.19 10 t − = × seconds, corresponding to the 64 th frame of the simulation. The position-domain particle reflects off the boundary while the momentum-domain particle travels through the boundary and reappears in the opposite boundary, travelling in the same direction at the same velocity. Figure 1 charts the relative position along the x-axis of the position-and momentum-domain particles. The position is given by the expectation value X . A marker has been placed at the 64 th frame to show where the paths are expected to diverge. Four free-particle animations are produced for the position domain. One overhead view and one crosssectional view are produced for the finite-difference free-particle simulation and one overhead view and one cross-sectional view are produced for the spectral free-particle simulation. Figure 2 shows the cross-sectional view of the free particle's components and density produced by finitedifference methods. The cross-section is along the x-axis in the position domain. Figure 3 shows the overhead view of the free particle's position-domain density produced by finite-difference methods. The particle diffracts as it reflects off the boundary. Figure 4 shows the cross-sectional view of the free particle's components and density produced by spectral methods. The cross-section is along the x-axis in the position domain. . Cross-sectional animation still at 64 th frame of free particle simulation produced by spectral methods. Figure 5 shows the overhead view of the free particle's position-domain density produced by spectral methods. The particle continues to diffuse in space as it pass through the boundary.
Square Well
Based on the initial conditions and predicted velocity, the particle is expected to reach the boundary at time . Comparing X of bound particle in square well produced by finitedifference and spectral methods.
Position Domain
Four square-well animations are produced for the position domain. One overhead view and one cross-sectional view are produced for the finite-difference square-well simulation and one overhead view and one crosssectional view are produced for the spectral square-well simulation. Figure 7 shows the cross-sectional view of the bound particle's components and density produced by finitedifference methods. The cross-section is along the x-axis in the position domain. Figure 8 shows the overhead view of the bound particle's position-domain density produced by finitedifference methods. The particle diffracts as it reflects off the potential wall. Figure 9 shows the cross-sectional view of the bound particle's components and density produced by spectral methods. The cross-section is along the x-axis in the position domain. Figure 10 shows the overhead view of the bound particle's position-domain density produced by spectral methods. The particle diffracts as it reflects off the potential wall.
Momentum Domain
One additional animation is produced for the momentum-domain, square-well simulation that shows the crosssectional view of the momentum density function. The cross-section is taken along the x k axis. Figure 11 shows a cross-sectional view of the density function in the momentum domain at the 42 nd frame. As the particle interacts with the electrostatic potential, the particle is reflected and reverses direction. In the momentum domain, this is indicated by the density function disappearing from the positive axis and reappearing on the negative axis. Figure 12 illustrates this reversal of direction over time by measuring the expectation value x K . A marker has been placed at the 42 nd frame indicating when the particle interacts with the boundary and reverses direction.
Discussion
Momentum-domain simulations of the time-dependent Schrödinger equation provide precise and accurate results; however, the application of these techniques is not limited to the Schrödinger equation. The methods described in this paper are also suitable to simulate the heat equation , where f describes temperature in space and time and α is the thermal diffusivity. The spectral methods described here may be applied to any parabolic differential equation. The spectral methods also transform the Hartree-Fock operation in a many-body problem from convolution in the position domain to entry-wise multiplication in the momentum domain, although this application is not explored here due to resource constraints.
These simulations were produced on single-core, AMD Athlon X2 processor. The spectral methods demonstrated faster performance for the free-particle simulation, while the finite difference methods demonstrated faster performance for the square-well simulation. None of the simulations employed parallel computing techniques due to the limitations of the hardware. Multiple cores would allow multiple Fourier transforms to be calculated at the same time. Because the spectral-method, square-well simulation requires multiple Fourier transforms at every time step, introducing parallel computing techniques would increase performance substantially. | 4,002.6 | 2015-08-20T00:00:00.000 | [
"Physics"
] |
Estimation of the ROC curve and the area under it with complex survey data
Logistic regression models are widely applied in daily practice. Hence, it is necessary to ensure they have an adequate predictive performance, which is usually estimated by means of the receiver operating characteristic (ROC) curve and the area under it (area under the curve [AUC]). Traditional estimators of these parameters are thought to be applied to simple random samples but are not appropriate for complex survey data. The goal of this work is to propose new weighted estimators for the ROC curve and AUC based on sampling weights which, in the context of complex survey data, indicate the number of units that each sampled observation represents in the population. The behaviour of the proposed estimators is evaluated and compared with the traditional unweighted ones by means of a simulation study. Finally, weighted and unweighted ROC curve and AUC estimators are applied to real survey data in order to compare the estimates in a real scenario. The results suggest the use of the weighted estimators proposed in this work in order to obtain unbiassed estimates for the ROC curve and AUC of logistic regression models fitted to complex survey data.
| INTRODUCTION
Prediction models are widely used in many different fields.Medical research, meteorology, business, and biology are just a few examples.
Although prediction models can be used for a variety of purposes, they are highly used for decision-making.For example, in finance, they can be useful for predicting loan defaults (Li et al., 2022); in ecology, particularly in fisheries, they are often used to make conservation decisions (Guisan et al., 2013;Li et al., 2020); medicine is another field in which prediction models are widely implemented as a support for decision-making, where, they can be helpful for deciding whether a patient should be admitted to an intensive care unit or not, among other purposes (Arostegui et al., 2019).Therefore, given the impact that the use of prediction models can have in many situations, it is necessary to ensure that these models are valid and applicable in practice.Thus, several aspects need to be considered in the development process of these models.A useful checklist can be found in (Steyerberg, 2009).In particular, when the goal is prediction, ensuring good model performance is essential.In this work, we focus on logistic regression models for dichotomous response variables.Model performance of logistic regression models is usually analysed by means of calibration and discrimination ability (Steyerberg, 2009).Calibration measures the agreement between outcomes and predictions (see, e.g., the goodness-of-fit test proposed by Hosmer and Lemesbow (1980)).In this study, we bring discrimination ability into focus, which measures the ability of the models to distinguish between units with the event of interest and without it.This is usually measured by means of the receiver operating characteristic (ROC) curve, which is defined as the curve formed by specificity and sensitivity parameters (i.e., probability of properly classifying individuals without and with the event of interest, respectively) across all the possible cutoff points (Green & Swets, 1966;Pepe, 2003;Swets & Pickett, 1982).The area under the ROC curve (AUC), is one of the most widely used summary measures to analyse the discrimination ability of logistic regression models (Pepe, 2003).Bamber (1975) showed the equivalence between the area under the ROC curve and the Mann-Whitney U-statistic, offering in this way an interesting interpretation of the AUC as the probability that an individual with the event of interest is given by the model a higher probability of event than an individual without the event of interest.
Complex survey data are becoming increasingly popular in various fields including health and social sciences, among others (see, e.g., Fisher et al., 2020).This type of data is collected from a finite population, concerned to be studied, by some complex sampling design, such as stratification, clustering, or a combination of them in different stages of the sampling process.One of the differences between complex survey data and simple random samples is that, in the first, each sampled observation ðiÞ has assigned a sampling weight w i , which is defined as the inverse probability weight or in other words the probability that unit i is included into the sample S, that is, w i ¼ 1=π i where π i ¼ Pði SÞ.The sampling weight assigned to each observation indicates the number of units that this observation represents in the finite population.Due to these particularities, the straightforward application of the most commonly applied statistical techniques for the development of prediction models, which assume that the data have been randomly collected and that sampled units are independent and identically distributed, is usually not appropriate for complex survey data and needs to be checked before being implemented in this context (Heeringa et al., 2017;Skinner et al., 1989).
For this reason, the effect of complex survey data on the development of prediction models in general, and logistic regression models in particular, is being widely discussed in the literature in recent years.One of the most discussed topics in the context of complex surveys is the effect of the sampling design on the estimation of model parameters (see, e.g., Binder & Roberts, 2009;Iparragirre et al., 2023;Lumley & Scott, 2017;Scott & Wild, 1986, 2002 as a summary of a large debate on this topic).Similarly, complex survey data have shown to have a great impact on the development of prediction models, and numerous advances are being made in the last years in this field.Lumley and Scott (2015) proposed new design-based estimators for estimating two widely used parameters for model selection, Akaike Information Criterion (AIC) and Bayesian Information Criterion (BIC), considering the sampling design.Iparragirre et al. (2023) proposed a new technique for considering complex sampling designs for variable selection with lasso regression models.Focusing on the evaluation of model performance, Archer et al. (2007) proposed a goodnessof-fit test that considers complex sampling designs to analyse the calibration of the models fitted to complex survey data.In the context of the discrimination ability, Yao et al. (2015) proposed a modification of the Mann-Whitney U-statistic in order to consider the sampling design to estimate the AUC of the models, incorporating pairwise sampling weights, which are defined as the inverse joint inclusion probability of a pair of observations ði, jÞ, that is, w ij ¼ 1=π ij where π ij ¼ P½ði SÞ \ ðj SÞ, 8i, j S (Horvitz & Thompson, 1952;Särndal et al., 2003).Finally, Iparragirre et al. (2022) proposed incorporating sampling weights into the estimation process of optimal cutoff points for individual classification as units with or without the event of interest.
As mentioned previously, in this work, we aim to focus on the evaluation of the discrimination ability of logistic regression models.Even though Yao et al. (2015) proposed a weighted estimator for the AUC, to our knowledge, there is a lack of proposals for estimating the ROC curve considering complex sampling designs.Therefore, the main goal of this work is to propose a weighted estimator for the ROC curve.In particular, we propose to use weighted specificity and sensitivity estimators defined in Iparragirre et al. (2022) to define a new weighted estimator for the ROC curve.In addition, we calculate the area under the curve in order to estimate the AUC parameter following Bamber (1975) and Tsuruta and Bax (2006), and finally, we show that this AUC estimator defined as the area under the weighted estimate of the ROC curve is equal to the weighted Mann-Whitney U-statistic considering marginal sampling weights w i , 8i S, rather than pairwise sampling weights w ij , 8i, j S as proposed by Yao et al. (2015).The estimation of the AUC is then a simple weighted expression that can easily be calculated in practice, given that the marginal sampling weights are usually explicitly available when working with complex survey data, in contrast to the pairwise sampling weights, which usually need to be calculated by means of some computational package.The performance of this proposal is analysed by means of a simulation study, in which the weighted and unweighted estimates of the ROC curve and AUC are compared with the true population ones.In addition, the proposed methods are applied to real survey data, and the weighted estimates of the ROC curve and AUC are compared with the unweighted ones.
The rest of the paper is organised as follows.In Section 2, we first set the basic notation.Then, we define the proposed weighted estimator of the ROC curve and we calculate the area under it.Finally, we show the equivalence between the area under the weighted estimate of the ROC curve and the weighted Mann-Whitney U-statistic considering marginal sampling weights.In Section 3, the simulation study conducted in order to analyse the performance of the proposed estimators is defined and the results obtained are depict and summarised.In Section 4, the proposed estimators are applied to real survey data.Finally, the paper concludes with a discussion in Section 5.
| METHODS
The goal of this section is to describe our proposal to estimate the ROC curve of logistic regression models fitted to complex survey data considering sampling weights.We calculate the area under the weighted estimate of the ROC curve in order to estimate the AUC, and we show the equivalence between this area and a modification of the Mann-Whitney U-statistic considering sampling weights, which leads us to conclude that this estimator can be used in order to obtain weighted estimates of the AUC.
The rest of the section is organised as follows.In Section 2.1 we denote the basic notation related to the logistic regression model, ROC curve, and AUC, as well as, complex survey data.In Section 2.2.1, we define our proposal to consider sampling weights to estimate the ROC curve and the area under the curve (AUC).In Section 2.2.2 we show the equivalence between the area under the weighted estimate of the ROC curve and the Mann-Whitney U-statistic considering sampling weights.
| Background and basic notation
Let X be the vector of covariates and Y the dichotomous response variable, which takes the value Y ¼ 1 for the units with the characteristic of interest (events), and Y ¼ 0 otherwise (non-events).Let the conditional probability of event for an individual i given the values of its vector of covariates x i .Let β indicate the vector of regression coefficients.Then the specific form of the logistic regression model is as follows: Based on the probability pðx i Þ and a cutoff point c, each individual can be classified as event (if However, this classification may be correct or incorrect depending on the selected cutoff point.The correct classifications, based on a particular cutoff point, are usually quantified by specificity (SpðcÞ) and sensitivity ðSeðcÞÞ parameters, which are defined as the probabilities of correctly classifying the nonevents and events, respectively, that is, The discrimination ability of a logistic regression model is usually evaluated by means of the area under the ROC curve (AUC), where the ROC curve is defined as the set of pairs 1 À SpðcÞ and SeðcÞ across all the possible cutoff points c (Green & Swets, 1966;Swets & Pickett, 1982).
The AUC ranges from 0.5 (an uninformative model) to 1 (a perfect model in terms of discrimination) (Steyerberg, 2009).
Let S indicate a sample of n observations of the vector of random variables ðY, XÞ, that is, Let S 0 and S 1 be its subsamples of sizes n 0 and n 1 formed by the units without the event of interest and with the event of interest, respectively (note that S 0 \ S 1 ¼ ; and S 0 [ S 1 ¼ S).
Let β indicate the vector of estimated regression coefficients, which are usually estimated by means of the likelihood function in Equation (3) and pi ¼ pðx i Þ the corresponding estimated probabilities of event, 8i S (McCullagh & Nelder, 1989): In practice, specificity and sensitivity parameters for a particular cutoff point c are estimated as proportions of correctly classified sampled non-events and events, respectively (see, e.g., Pepe;2003), that is, Iðp j < cÞ and c SeðcÞ ¼ where IðÁÞ denotes the indicator function.Then, the estimated ROC curve is defined by means of each estimated pair of sensitivity and specificity parameters, for each possible cutoff point (Pepe, 2003) as shown in Equation ( 5): Bamber (1975) showed that the area under the ROC curve defined in Equation ( 5) can be estimated ( d AUC) as described in Equation ( 6), by means of the Mann-Whitney U-statistic: However, in the context of complex survey data, sample S is usually obtained by sampling a finite population U ¼ f1,…, Ng of interest for the survey, following some complex sampling design.For each sampled observation 8i S, w i ¼ 1=π i , ðwhereπ i ¼ Pði SÞÞ denotes the sampling weight assigned to this observation, indicating the number of units from the finite population it represents.In this context, the regression coefficients and the corresponding probabilities of events are usually estimated ( β and pi ¼ pðx i Þ, 8i S) maximising the pseudo-likelihood function (Binder, 1983;Iparragirre et al., 2023) defined in Equation ( 7): In practice, the goal is to evaluate the fitted model's performance in the finite population U, that is, the ability of the model fitted following Equation ( 7) to discriminate individuals with and without the event of interest in the finite population.Let ðy i , x i Þ f g N i¼1 be N realisations of the set of random variables ðY, XÞ, and let N 0 and N 1 indicate the size of the subsets formed by the nonevents (U 0 ) and events (U 1 ) of the finite population U.Then, the finite population ROC curve can be defined as in Equation ( 5) where sensitivity and specificity parameters are estimated following Equation ( 4) for all the units in the finite population U, that is, The fitted model's AUC in the finite population could then be defined as follows: However, note that the set of ðy i , x i Þ f g N i¼1 is usually unknown, except for the sampled units i S, so the finite population ROC curve and AUC need to be estimated based uniquely on S. We believe that in the context of complex survey data, if the ROC curve and the AUC of the fitted model are estimated based on Equations ( 5) and ( 6), which were designed to be applied in simple random samples and do not consider the sampling weights, then biassed estimates can be obtained.For this reason, we propose a new estimator for the ROC curve and the AUC which considers the sampling weights.This proposal is described in Section 2.2 below.
| Proposal
In Section 2.2.1, we propose an estimator to estimate the ROC curve for logistic regression models fitted with complex survey data and the AUC as the area under the curve.In Section 2.2.2, we show the equivalence between the proposed AUC estimator and the Mann-Whitney U-statistic incorporating marginal sampling weights.
| Estimation of the ROC curve and the area under it
We propose to estimate the ROC curve considering the sampling weights, as follows: for which specificity and sensitivity parameters are estimated by means of the sampling weights (Iparragirre et al., 2022): and c Se w ðcÞ ¼ Therefore, we propose to calculate the area under d ROC w ðÁÞ in order to estimate the AUC (Tsuruta & Bax, 2006).Let us denote as A the area under the curve.We now proceed to describe how the area under the ROC curve defined in Equation ( 10) can be calculated.Note that in practice, we always work with finite sample sizes, and hence, the number of different estimated probabilities is finite.Let us denote as q the total number of different estimated probabilities, that is, pðqÞ < … < pð1Þ (where q ≤ n, being q ¼ n if and only if all the estimated probabilities for each sampled unit are different).Note that for every cutoff point chosen between two ordered probabilities, the same values for the specificity and sensitivity parameters will be obtained, and therefore, the same point would be defined in the ROC curve.Then, the ROC curve will be completely defined with q þ 1 different cutoff points.Specifically, the smallest possible cutoff point is c q < pðqÞ , which will classify all the sampled units as events and therefore the estimate of the sensitivity will be 1 and the specificity will be 0 (see Equation ( 11)), that is, the cutoff point c q will draw the point ð1 À c Sp w ðc q Þ, c Se w ðc q ÞÞ ¼ ð1,1Þ in the ROC curve.In the same way, the point drawn in the ROC curve for c 0 > pð1Þ will be the point Se w ðc 0 ÞÞ ¼ ð0,0Þ.Let us denote and sort the rest of the q À 1 cutoff points as follows: For ease of notation, 8m ¼ 1, …, q À 1, each cutoff point c m can be defined as the average value of the probabilities pðmþ1Þ and pðmÞ , that is, Note that in this way, all the defined cutoff points will be different from the estimated probabilities and since between any two different ordered predicted probabilities, a cutoff point has been defined, only one different predicted probability lies in the interval Se w ðc m ÞÞ.In this way, the estimated ROC curve will be a polygonal line defined by q segments.Each of these segments will define an area with the abscissa axis.Let us denote as A m , 8m f1,…, qg each of these areas.A graphical explanation can be seen in Figure 1.
We now proceed to calculate analytically the area under the ROC curve defined in Equation ( 10).In particular, as the area A 1 is a triangle of base 1 À c Sp w ðc 1 Þ h i and height c Se w ðc 1 Þ, it can be calculated as follows: F I G U R E 1 Graphical explanation of the empirical weighted receiver operating characteristic (ROC) curve, where For m ¼ 2,…, q, the areas A m are right-angled trapezoids, the area of which can be easily calculated as the sum of the triangle Se w ðc mÀ1 Þ and rectangle A 2 m of the same base and height c Se w ðc mÀ1 Þ: Then, the area under the ROC w ðÁÞ curve ðAÞ can be calculated as the sum of the areas defined in Equations ( 14) and ( 15).Note that c Se w ðc 0 Þ ¼ 0 and c Sp w ðc 0 Þ ¼ 1.Then, Equation ( 14) that defines A 1 can be rewritten in terms of those values for convenience.Finally, the area under the curve can be easily calculated as follows: 2.2.2 | Equivalence between the area under the d ROC w ðÁÞ curve and Mann-Whitney U-statistic We propose to incorporate the marginal sampling weights into the Mann-Whitney U-Statistic as follows to estimate the weighted AUC: In the following lines, we show that the area under the estimated ROC curve defined in Equation ( 16) is equivalent to the Mann-Whitney U-statistic considering marginal sampling weights as defined in Equation ( 17).In order to prove the equivalence between both approaches, our goal is to rewrite Equation ( 17) in terms of sensitivity and specificity parameters.Let us rewrite it as follows as the first step: Then, we can rewrite the expressions Iðp j < pk Þ and Iðp j ¼ pk Þ as a function of the previously defined cutoff points.Given that Then, 8j S 0 the inequality pj < pk will be satisfied if and only if pj < c m , as graphically shown in Figure 2. Thus, note that Iðp j < pk Þ can be rewritten as follows: Then, following Equation ( 19) and the definitions given in Equation ( 11), let us rewrite the first term of Equation ( 18) in terms of sensitivity and specificity parameters as follows: In the same way, we will now proceed to rewrite the expression Iðp j ¼ pk Þ.As stated above, 8k S 1 , 9! m f1,…, qg : pk c m , c mÀ1 ½ Þ .Thus, 8j S 0 , the equality pj ¼ pk will only be satisfied if pj is in the same range as pk , that is, pj c m , c mÀ1 ½ Þ (see Figure 3).Let us rewrite Following Equation ( 21) and the definitions given in Equation ( 11), let us rewrite the second therm of Equation ( 18) in terms of sensitivity and specificity parameters as follows: This image is intended to be helpful to better understand Equation ( 19) and indicates in which situations F I G U R E 3 This image is intended to be helpful to better understand Equation ( 21) and indicates in which situations Finally, Equation ( 18) can be rewritten as the sum of Equations ( 20) and ( 22): Note that Equations ( 16) and ( 23) are equal, so we have shown that A ¼ d AUC w .
| SIMULATION STUDY
The goal of this simulation study is to analyse the performance of the proposed estimators in comparison with the traditional unweighted estimators of the ROC curve and AUC.In Section 3.1, we describe the data generation process and different scenarios considered throughout the study; in Section 3.2, we describe the simulation study process, and finally, in Section 3.3, we summarise the main results.
| Data generation and scenarios
In the following lines, the data simulation process is described.Let us define as N ¼ 10,000 the finite population size.A set of p ¼ 5 covariates (X 1 , …,X 5 ) and two latent variables (Z 1 and Z 2 , which are used to define the response variable and the sampling design, but are not available in the samples when fitting models) have been generated.
A total of three different scenarios have been defined based on different sampling designs.On the one hand, a stratified sampling design without clustering was defined (let us denote this scenario as SH, hereinafter), in which different strata are defined in the finite population, and a number of individuals are sampled from each stratum.On the other hand, we defined a stratified sampling design with clustering (scenario SC), in which different strata are defined in the finite population, a number of clusters or groups of units are selected from each stratum, and finally, a number of individuals are sampled from each selected cluster.In addition, in this scenario, SC two situations have been distinguished: First, all the variables have been considered as unit-level variables (we denote this scenario as SC.0, given that there are d ¼ 0 cluster-level variables); and, second, in the other scenario, one cluster-level variable (d ¼ 1) has been considered (scenario SC.1).Note that in scenario SH, all the variables must be defined at unit-level (d ¼ 0) since there is no cluster.We proceed below to explain the data generation process for each of these scenarios: 1.For d ¼ 0 (SH and SC.0) and d ¼ 1 (SC.1),N realisations have been made following the Gaussian distribution defined in Equation ( 24): where μ ðpÀdÞ indicates the null vector of dimension 1 Â ðp À 1Þ and Σ ðpÀdÞÂðpÀdÞ a matrix of dimension ðp À dÞ Â ðp À dÞ defined by values of 1 on the diagonal and η ¼ 0:15 off-diagonal, i.e., μ ðpÀdÞ ¼ ð0,…,0Þ T and Σ ðpÀdÞÂðpÀdÞ ¼ ð1 À ηÞ being I ðpÀdÞÂðpÀdÞ the identity matrix and J ðpÀdÞÂðpÀdÞ the matrix of 1s.
2. Let us denote as fz i ¼ ðz i,1 , z i,2 Þg N i¼1 the set of N realisations of Z 1 and Z 2 .Data are sort based on z i β Z , 8i ¼ 1,…,N, where Strata are defined by partitioning the ordered population data set on sets of the same size (H ¼ 10 strata) in all the scenarios, being each stratum of size N h ¼ 1000, 8h ¼ 1,…, H.In addition, in scenarios SC.0 and SC.1, each stratum has been partitioned into A h ¼ 10 clusters 8h ¼ 1,…, H.In this way, a total of A * ¼ 100 clusters of size N h,α ¼ 100 are generated, 8h ¼ 1, …, H and 8α ¼ 1,…A h .
3. If d ¼ 1, then X 1 is a cluster-level variable (SC.1).We generate it by making A * ¼ P H h¼1 A h realisations of X 1 $ Nð0,1Þ.Note that for two different units in the same cluster, their corresponding cluster-level covariates should take the same value, that is, 8i, j in the same cluster, x i,1 ¼ x j,1 .Therefore, we repeat each realisation N h,α times.4. We now have defined the values corresponding to X 1 ,…, X 5 variables for all the units in the finite population: fx i ¼ ðx i,1 , …, x i,5 Þg N i¼1 .Let us define β X as follows: Then, we generate the probabilities of event as follows: and the value for the response variable y i is randomly generated by following Bernoulliðpðx i , z i ÞÞ.We set β 0 ¼ À5, defining in this way a prevalence (i.e., probability of event) of around 25%.
The finite population U is defined as the set of values corresponding to the response variable y i and the covariates x i (excluding the latent variables z i ), 8i ¼ 1, …, N as well as strata and cluster indicators corresponding to each of them.
5. Different sampling schemes have been considered in this simulation study.On the one hand, in the scenario in which a stratified sampling design without clustering is defined (SH), the following number of units have been considered from each stratum (n h , 8h ¼ 1, …H): On the other hand, in the scenarios in which a stratified sampling design with clustering is considered (SC.0 and SC.1), a h ¼ 2, 8h ¼ 1,…, H clusters have been sampled in the first place.Then, from each sampled cluster of stratum h, the following number of units have been sampled (n h,α , 8α ¼ 1,2): It should be noted that due to the way in which the sample design has been defined, the probabilities of event given the covariates are roughly ordered from highest to lowest in the different strata.Therefore, by sampling many units from the strata at the edges, we are sampling more individuals with higher and lower probabilities (scheme (a)).In contrast, when sampling more individuals from the central strata, more individuals with medium probabilities of event are sampled (scheme (b)).The two sampling schemes differ on this point.
6. Depending on the sampling design defined in each scenario, sampling weights are calculated as follows.In scenarios without clustering (SH), the sampling weight for each unit i from stratum h is defined as the total number of units in the population in stratum h (N h ) divided by the number of units sampled from it (n h ), that is, In the same way, in the scenarios which consider clustering (SC.0 and SC.1), the sampling weight for each sampled unit i S from cluster α from stratum h is defined as follows:
| Setup
Considering the scenarios described in Section 3.1, a finite population was simulated in each scenario.The theoretical model is fitted to the finite population, and the ROC curve and AUC of this model are calculated following Equations ( 8) and ( 9) (let us denote as ROC theo and AUC theo the theoretical ROC curve and AUC, corresponding to the finite population model).Note that these parameters measure the performance of the theoretical finite population model.Each population is sampled R ¼ 500 times, following in each case the corresponding complex sampling design.In each of the samples, a weighted logistic regression model was fitted and its ROC curve and AUC were estimated, ignoring the sampling weights (unweighted method) and considering them (weighted method).Note that in practice, we aim to analyse how those estimators perform in order to estimate the fitted model's ROC curve and AUC in the finite population.Therefore, in order to analyse and compare the performance of both estimators, we compare each of the estimates to the true finite population ROC curve and AUC estimates of the model fitted to the sample (which is calculated by extending the fitted sample model to the finite population), rather than to the theoretical population model parameters.
We will denote these true finite population parameters as d ROC pop and d AUC pop .Note that these parameters indicate the true performance of the fitted sample model in the finite population.This process is described in detail below and summarised in Figure 4.For r ¼ 1, …, R: Step 1. Obtain a sample S r & U by means of one of the sampling schemes described in Section 3.1 and calculate the sampling weights w r i , 8i S r following the corresponding Equation ( 29) or (30).
Step 2. Fit the model to S r maximising the pseudo-likelihood function in Equation ( 7 17)), to obtain unweighted and weighted estimates, respectively.In addition, we estimate the AUC by means of pairwise sampling weights following the proposal of Yao et al. (2015), which will be denoted as d AUC r pairw .
Step 4. By means of the βr estimated in Step 2., estimate the probabilities of event for all the units in the finite population, pr i , 8i ¼ 1, …, N. Estimate the true ROC curve and AUC in the population following Equations ( 8) and ( 9 F I G U R E 4 Graphical explanation of the simulation study setup. Step 5. Calculate the difference between the unweighted or weighted estimates and the true population AUC (obtained based on the model fitted to the sample): In addition, in order to compare our proposal that considers marginal sampling weights to the proposal considering pairwise sampling weights, we define the difference between pairwise estimates to the true population model and to our proposal as follows: All computations were performed in (64 bit) R 4.2.2 (R Core Team, 2022) and a MacBook Pro equipped with 16 GB of RAM, Apple M1 Chip, and macOS Monterey 12.2.1 operating system.Logistic regression models were fitted by means of survey R package Lumley (2010Lumley ( , 2020)).
| Results
In this section, we summarise the main results we obtained from the simulation study.In Figure 5, the theoretical ROC curve of the finite population model, as well as the true population ROC curves and the weighted and unweighted estimates obtained based on the models fitted across R ¼ 500 samples are shown.Figure 6 depicts the boxplots of the differences between the unweighted (diff r unw ) and weighted (diff r w ) estimates and the true population AUC of the models fitted to the samples (see Equation ( 31)), while Figure 7 depicts the boxplots of the differences between the AUC estimates obtained by means of the pairwise and marginal sampling weights (wdiff r ).Table 1 summarises the numerical results.Due to F I G U R E 5 Unweighted (unw; see Equation ( 5)) and weighted (w, Equation ( 10)) estimates of the receiver operating characteristic curves, as well as the true population ROC curves (pop) of the models fitted across r ¼ 1,…, 500 samples, together with the theoretical ROC curve (theo) of the model fitted to the finite population in each scenario drawn in the simulation study.
the large number of results we obtained, we begin by summarising the main conclusions, and then we proceed to analyse the differences between the different scenarios.
As shown in Figure 5, the theoretical ROC curve of the population model is above most of the population ROC curves of the models fitted to the samples.Similarly, as can be regarded in Table 1, the average true population AUCs are lower than the theoretical AUCs.This indicates that population models have a greater discrimination ability than the models fitted to the samples.Therefore, in order to make fairer comparisons and compare the AUCs of the same models, we compare the ROC curve and AUC estimates obtained with different methods, to the true population parameters rather than the theoretical ones.In general terms, the results of the simulation study show that, under the scenarios that have been considered, the weighted estimates of the AUC are closer than the unweighted ones to the true population AUC.The weighted estimates are slightly optimistic on average, given that a bit greater AUCs than the true ones have been estimated.In contrast, unweighted estimates sometimes overestimate the true finite population AUC, and other times underestimate it, depending on the scenario (in any case, showing a greater absolute bias than the weighted estimates).In terms of variability, no major differences have been observed between the two estimators, and, depending on the scenario, one estimator or the other shows more variability.The marginal and pairwise weighted estimators perform quite similarly in all the scenarios, both in terms of bias as well as variability.However, it should be noted that as shown in Figure 7, the estimates based on pairwise sampling weights are slightly greater than the ones obtained based on marginal sampling weights.Thus, the estimates based on pairwise weights overestimate a little more the true population AUC than the estimates based on marginal weights, even though, those differences are minimal in terms of bias.In contrast, computation times are considerably improved with the estimator proposed in this work (up to five times more efficient as can be seen in Table 1), given that the pairwise sampling weights need to be calculated for each particular sampled pair, in contrast to the marginal ones, which are easily available in most cases when working with this kind of data.
We now proceed to comment on the results in more detail.The sampling schemes (a) and (b) differ in the number of units sampled from each stratum, and more specifically, in the number of units sampled with (a) higher and lower predicted probabilities or (b) central predicted probabilities.These differences have an effect on the unweighted estimates in terms of the difference in comparison with the true population AUC, which in scenarios with sampling scheme (a) underestimates and in scenarios with sampling scheme (b) overestimates it.In contrast, for the weighted estimates, no great differences have been observed in terms of difference in comparison to the true population AUC depending on the sampling schemes (a) and (b).For example, as can be observed in Table 1, in Scenario SC.0 (a) the average difference between the unweighted estimates F I G U R E 6 Boxplots of the difference (see Equation ( 31)) between the estimated area under the curves (AUCs) by means of the unweighted (unw, Equation ( 6)) and weighted (w, Equation ( 17)) estimators and the true population AUC of the models fitted across r ¼ 1,…, 500 samples in all the scenarios drawn in the simulation study.and the true population AUC is À0.081, while in Scenario SC.0 (b) the average difference is 0.073.For the weighted estimates, under the same scenarios, the average differences are 0.005 and 0.008, respectively.These differences can also be observed in Figure 5, where the unweighted ROC curves are under the true population ROC curves in scenarios (a), while in scenarios (b) the unweighted ROC curves are over the true ones, as well as, over the weighted ones, indicating that the unweighted estimates overestimate more than the weighted ones in these scenarios.
However, in terms of variability, the performance of the unweighted and weighted estimates differ under sampling schemes (a) and (b).In scenarios considering sampling scheme (a), the variability of the unweighted estimates is greater than the variability of the weighted ones, while in scenarios considering sampling scheme (b), the difference is reversed.As shown in Table 1, in Scenario SH (a), the standard deviation of the unweighted estimates is 0.018, slightly greater than the variability of the weighted estimates which is 0.014.In contrast, in Scenarios SH (b), the standard deviation of the unweighted and weighted estimates are 0.012 and 0.020, respectively.In addition, the variability of the unweighted estimates is greater in (a) than in (b) (for the weighted estimates this difference is not as remarkable as for the unweighted estimates).For example, in SC.0 (a) the standard deviation of the unweighted estimates (0.035) is 2.5 times greater than the standard deviation in SC.0 (b) (0.014).
Results also show that the performance of the two estimators differs depending on the sampling design.In particular, a greater optimism of the weighted estimates has been observed in scenarios with cluster-level variables SC.1 than in scenarios SC.0 and SH.For example, in scenario SC.1 (a), the average difference between the weighted estimates and the true population has been 0.023, while in scenario SH (a), the average difference has been 0.005.This effect can also be observed in Figure 6.The ROC curves depicted in Figure 5 also show that in Scenarios SC.1 (a) and SC.1 (b), most of the weighted ROC curves are above the true population curves, while in the rest of the scenarios, the true population ROC curves are more or less in the center of the weighted ROC curves' band.This effect has not been observed for the unweighted estimates.In contrast, the sampling design has affected the variability of both, unweighted and weighted estimates.Specifically, the standard deviation of the estimates in scenarios SH is lower than that in scenarios SC.0 which, in turn, is lower than the standard deviation in scenarios SC.1 (see Table 1 for more details).It should also be noted that the standard deviation of the true population AUCs across R ¼ 500 samples is greater in scenarios SC.1 than in the rest of the scenarios (Table 1).This can also be observed in Figure 5, where the true population ROC curves show the greatest variability in scenarios SC.1.
F I G U R E 7 Boxplots of the differences between the estimated area under the curves (AUCs) by means of the AUC estimator based on pairwise sampling weights and the one that considers marginal sampling weights (wdiff; see Equation ( 32)), when estimating the AUC of the models fitted across r ¼ 1,…, 500 samples in all the scenarios drawn in the simulation study.
The methodology proposed in Section 2 has been applied to the Survey on the Information Society in Companies 1 (ESIE survey), which was described in detail in Iparragirre et al. (2022).This survey was carried out among the companies in the Basque Country (BC) in order to collect information about the use of technology.In particular, the response variable considered here is the same as the one used in the above-mentioned study, that is, a dichotomous response variable indicating whether a company has its own web page (1) or not (0).A sample of n ¼ 7725 was considered for this application, and the AUC of the model fitted in Iparragirre et al. (2023) was estimated.Covariates included in the model represent the activity of the company, the number of employees, and the ownership.
The unweighted and weighted AUC estimates and the corresponding Bootstrap 95% confidence intervals (CIs) are shown in Table 2.For the 95% CI of the unweighted estimate, the Bootstrap CI is calculated by means of the pROC R package (Robin et al., 2011), while the 95% CI of the weighted estimate is calculated by generating Bootstrap resamples based on replicate weights (Rao & Wu, 1988) using the survey R package (Lumley, 2020), both of them considering B ¼ 2000 Bootstrap resamples.The unweighted and weighted ROC curve estimates are depicted in Figure 8.Note that in this case, as we are working with real survey data; we cannot know which the true population ROC curve and AUC are.
Even though the differences between the unweighted and weighted estimates are not as large as the ones analysed in the simulation study, the unweighted estimate is larger than the weighted estimate, as it happens in scenarios (b) of the simulation study.Considering the results of the simulation study, we can assume that the weighted estimate will be a bit above the true population AUC, and therefore, we can conclude that probably the unweighted estimate of the AUC is overestimating it.In addition, note that the overlap between the two CIs is very slight.
T A B L E 1 Numerical results of the minimum value (min), maximum value (max), average (mean) and standard deviation (sd) of the population AUC (pop, Equation ( 9)), unweighted (unw, Equation ( 6)), weighted (w, Equation ( 17)) and pairwise (pairw) (Yao et al., 2015) estimates of the AUC of the models fitted across r ¼ 1,…,500 samples.Note: The average difference (av.diff) of the unweighted, weighted and pairwise estimates to the true population AUC estimates (see Equations ( 31) and ( 32)) with their standard deviations (sd), and the average computational times (av.time, in seconds) of each method with their standard deviations (sd) are also shown.In addition, the theoretical AUC (AUC theo ) of the finite population model in each scenario is available.Abbreviation: AUC, area area under the ROC curve.
In this work, we propose two new weighted estimators to estimate the ROC curve and AUC of logistic regression models fitted to complex survey data.In addition, we show that the area under the proposed weighted estimator of the ROC curve is equivalent to the weighted Mann-Whitney U-statistic incorporating marginal sampling weights, which are defined as the inverse probability weights for each sampled unit.A simulation study has been conducted in order to analyse the performance of the proposed estimators, and they have also been applied to real survey data.
The results of the simulation study suggest the use of the proposed weighted estimators rather than the unweighted ones.The unweighted estimators overestimate or underestimate the true population parameters, depending on the proportion of units sampled from each stratum.In particular, as more units with extreme (higher and lower) predicted probabilities are sampled in proportion, more nonevents with higher predicted probabilities as well as events with lower predicted probabilities are also sampled.This results in a lower estimate of the AUC.In contrast, as more central predicted probabilities are sampled, less extreme (higher and lower) predicted probabilities than necessary to properly represent the finite population will be sampled, leading to a greater estimate of the AUC due to the same reason.Weighted estimates correct for this bias providing ROC curve and AUC estimates that are closer to the true finite population parameters, since the presence of sampling weights gives each pair of individuals with and without the event of interest the relevance that they should have in representing the finite population.
F I G U R E 8 Weighted and unweighted receiver operating characteristic (ROC) curves of the models fitted to the ESIE survey data.
) by means of the covariate values x i and the sampling weights w r i , 8i S r (note that the latent variable values z i are only considered to define the sampling design and are not considered in the model estimation process).Obtain βr and the estimated probabilities of event pr i , 8i S r .Step 3. Estimate the ROC curve ( T A B L E 2 Estimated unweighted and weighted AUCs and the corresponding Bootstrap 95% CI of the model fitted to ESIE survey data. | 10,071.6 | 2023-01-01T00:00:00.000 | [
"Mathematics"
] |
A Comparison of RFID Anti-Collision Protocols for Tag Identification
: Radio Frequency Identification (RFID) is a technology that uses radio frequency signals to identify objects. RFID is one of the key technologies used by the Internet of Things (IoT). This technology enables communication between the main devices used in RFID, the reader and the tags. The tags share a communication channel. Therefore, if several tags attempt to send information at the same time, the reader will be unable to distinguish these signals. This is called the tag collision problem. This results in an increased time necessary for system identification and energy consumption. To minimize tag collisions, RFID readers must use an anti-collision protocol. Different types of anti-collision protocols have been proposed in the literature in order to solve this problem. This paper provides an update including some of the most relevant anti-collision protocols.
Introduction
The Internet of Things (IoT) is a set of millions of physical devices around the world that are connected to the Internet, collecting and sending data. It enables objects and machines to connect to the Internet and share collected information. Thanks to very cheap procedures and wireless networks, it is possible to turn anything, from a chip to a big building, into part of the IoT. This is an additional level of digital intelligence of devices that would otherwise be dumb, permitting communication without humans, and merging the digital and physical worlds. The term IoT was first coined by British entrepreneur Kevin Ashton in 1999, while working at Auto-ID Labs, referring specifically to a global network of objects connected by RFID [1]. Nowadays, the IoT has been the focus of numerous research studies. Moreover, although IoT extends beyond the RFID, this technology is one of the key cores to implement the IoT.
RFID technology uses radio frequency in order to identify tags, which are small sticks attached to the objects that wish to be identified [2,3]. The RFID system operates with a reader sending and receiving information from numerous tags in the interrogation range of its antenna at the same time. That is, the reader broadcasts messages using electromagnetic waves to tags in the interrogation area, while receiving responses from the tags through backscattering. Therefore, a multi-access arbitration procedure is necessary in order to prevent the response from being garbled [2,4].
When more than one tag transmits simultaneously to a reader, their backscattered signals cancel each other out, resulting in an illegible message by the reader. This is called collision and it causes a loss of identification time and an increase in the power consumed by the reader.
-
Signal interference occurs when the fields of two or more readers overlap and interfere. This problem can be solved by programming all readers to read at fractionally different times.
-
Multiple reads of the same tag occur when the same tag is read once by every overlapping reader.
• A tag collision occurs when more than one tag attempts to transmit its ID at the same time: the reader will receive a mixture of the tags' signals and cannot understand it. This type of collision is shown in Figure 2.
Simultaneous responses from numerous tags prevent the reader from correctly translating the signal, which decreases throughput. No tag is aware of the activity of any other tag, and so they cannot prevent the simultaneous transmission of tags. The transmissions of three tags, shown in Figure 1, are not synchronized, but in many cases, the reader is synchronized with at least one tag in the interrogation zone. In the presented illustration, the reader is prevented from decoding the entire transmission and it has experienced a collision. To solve this problem, anti-collision protocols are very influential.
Multi-Access Methods
Each anti-collision protocol uses certain multi-access methods for identification in order to physically separate the transmitters' signals. Accordingly, they can be categorized into four different types: Space Division Multiple Access (SDMA), Frequency Division Multiple Access (FDMA), Code Division Multiple Access (CDMA) and Time Division Multiple Access (TDMA) [2,5,8]. Figure 3 shows various multiple access and anti-collision procedures.
SDMA-The term space division multiple access relates to dividing of the channel capacity into separate areas. Protocols based on this method can point the beam at different areas in order to identify tags. The channel is spatially separated using complex directional antennas. Another means of achieving this is through the use of multiple readers. As a result, the channel capacity of adjoining readers is enhanced. A huge number of tags can be read simultaneously as a result of the spatial distribution over the entire layout. This method is quite expensive and requires complex antenna design. The use of this type of method is restricted to a few specialized applications [9][10][11]. This technique is shown in Figure 4. FDMA-Tags transmitting in one of several different frequency channels requiring a complex receiver at the reader. Consequently, different frequency ranges can be used for communication from and to the tags: from the reader to the tags, 135 kHz, and from the tags to the reader, in the 433-435 MHz range. However, this technique is expensive and is only intended for certain specific applications [12]. Figure 5 shows FDMA procedure. CDMA-Requires tags to multiply their ID by a pseudo-random sequence (PN) before transmission. CDMA is quite good in many ways, such as the security of the communications between the RFID tags and the reader, and multiple tag identification. It adds great complexity and is expensive for RFID tags. Furthermore, this method consumes a great deal of power and can be classified as a group with elevated demands [13]. Figure 6 shows this procedure.
TDMA-Given that it is less expensive, this is the most widely used method. This method involves the largest group of anti-collision algorithms. The transmission channel is divided between the participants and ensures that the reader can identify a tag at different times in order to avoid interfering with another one. The space-distributing characteristic of tags is not considered. The number of tags in the interrogation zone is reduced after every successful response. Another option involves the ability to mute all tags except for the transmitting tag. After that, the tags are activated one by one [14,15]. TDMA is shown in Figure 7. In an RFID environment, anti-collision protocols typically use the TDMA method. Protocols that use this method first select an individual tag from a large group using a specific algorithm and then the communication takes place between the selected tag and the reader. Significant increases in number of collisions in the identification process decreases the throughput and increases the number of transmitted bits. These protocols can be divided into three categories: Aloha-based protocols, tree-based protocols and hybrid protocols (which use a combination of the first two methods). In the following subsections, some Aloha, tree-based and hybrid protocols shall be presented.
Aloha Protocols
Aloha-based protocols use a random-access strategy in order to successfully identify the number of tags in an interrogation area [16][17][18][19][20]. They belong to the group of probabilistic protocols because the tags transmit their own ID in randomly selected slots in a frame in order to reduce the possibility of a collision. However, there is no guarantee that all of the tags will be identified in the interrogation process. These protocols suffer from the well-known tag starvation problem, in the sense that a tag may not be correctly read during a reading cycle due to an excessive number of collisions with that same tag. Every frame consists of a certain number of slots, and the tags can only respond once per frame [16].
The main Aloha-based protocols can be divided into four subgroups: Pure Aloha (PA), Slotted Aloha (SA), Frame Slotted Aloha (FSA) and Dynamic Frame Slotted Aloha (DFSA) protocols.
Pure Aloha
Pure Aloha (PA) is one of the simplest anti-collision protocols. It is based on TDMA [21,22]. Whenever tags enter the interrogation zone, they randomly choose a frequency on which to transmit their data. A collision will occur if several tags transmit data at the same time, resulting in complete or incomplete collisions. A complete collision occurs when the messages of two tags fully collide; an incomplete collision, however, takes place when only part of the tag message collides with another tag message. This procedure is shown in Figure 8 and will be repeated until all tags are successfully identified. PA has been presented using different extra-features [8,23] such as: muting for silencing tags after being identified; the 'Slow down' for decreasing a tag response rate after identification; the 'Fast Mode' for sending a silence message before a tag begins transmission; and combinations of these different features.
Slotted Aloha
To avoid incomplete collisions, Slotted Aloha (SA) has been created. In SA, the time is divided into several slots and each tag must randomly select a slot in which it will transmit its data [8,21,23,24]. The communication between the reader and the tag is now synchronous. An example of communication with this protocol is presented in Figure 9.
Also, SA can use features similar to those presented for PA. The muting of slow down features are used to silence or decrease the rate of the tags; the early end closes the slot earlier than normal; and also combinations of the types.
Framed Slotted Aloha and Dynamic Framed Slotted Aloha
In Framed Slotted Aloha (FSA), the time is divided into a variable number of frames and each frame consists of several slots [16,[25][26][27][28]. All tags need to transmit data into a fixed length frame, but each tag must choose only one slot in a frame to transmit data. This protocol significantly reduces the probability of collision since tags can only respond once in a frame. Before the onset of communication, the reader generates a random time that is less than the fixed frame size, to select just one slot in a frame. If a collision occurs, the involved tags will once again choose a slot in which to respond in the next frame. The main inconvenience of FSA is slot wastage when the number of tags is small and the size of the frame is significantly larger [29,30].
To ameliorate this disadvantage, the Dynamic Frame Slotted Aloha (DFSA) protocol was developed [16,31,32]. DFSA is capable of changing frame size according to an estimate of the number of tags. At the beginning of each frame, the reader informs the tags of the frame length. Every tag selects a random number [0, F − 1], where F denotes the frame size and all tags respond within the number of slots. At the end of the frame, the reader estimates the number of colliding tags, then adjusts F accordingly. There are certain disadvantages to tag estimation, such as: increased computational costs in the identification process and errors that degrade the protocol's efficiency. Examples of these two protocols are shown in Figure 10.
Q Protocol
The Q protocol is used in the EPCglobal Generation-2 (Gen-2) standard [18,33]. The Q protocol is a DFSA-type protocol that modifies frame size using feedback from the last frame accomplished [34]. The Q algorithm can jump into the following frame without finishing the current one [34].
The Q algorithm operates with two basic parameters: Q, and a constant c that can be modified depending on the situation [24,25]. The Q. variable is an integer ranging from 0 to 15. This protocol works with three types of commands: • Query command, transmitted by the reader to all tags in the interrogation area, in order to force all tags to choose a slot number (SN) from [0, 2 Q − 1]. This command initiates the identification process by providing a new value of Q.
•
QueryAdjust is the command used to instruct all tags to increase, decrease, or maintain the Q value unchanged and to reselect their SN. Q new denotes the last calculated Q. Accordingly, Q could be increased by c, decreased by c, or left unchanged, according to the algorithm.
•
QueryRep is used in order to notify all tags to decrease their SN by 1.
The procedure of the Q protocol appears in the flow chart in Figure 11. If it is time to initiate a new inventory round, the reader will transmit a Query command. If the tags receive Query or QueryAdjust, they need to choose SN from [0, 2 Q − 1]. If they receive a QueryRep, all unidentified tags decrease their SN counter by 1. Only tags with SN = 0 will generate a 16-bit random number (RN16) and will respond with an RN16 to the reader. There are three possibilities, depending on the tags' response: • Successful reply. If only one tag responds and the reader successfully received RN16. Subsequently, the reader sends the ACK and only the tag that successfully responded recognizes the ACK and reports its EPC to the reader.
•
Collided reply. If more than one tag transmits RN16, a collision occurs. Then, the tags will increase Q new by the constant c. Typical values for c are 0.1<c<0.5. The value of c is adjusted according to the type of application. Higher values of c will provide more aggressive frame adjustments.
•
No reply. When no tags respond in the slot, the reader decreases Q new by c.
A similar protocol solution provides different updating steps for no-reply and collided reply by calculating the probabilities of idle slots and collided slots, without considering the parameters in Gen-2 [28,35]. Table 1 summarizes the observations regarding the Aloha protocols, where the efficiency measures the exploitation of the tags' responses and the influence of a collision on a tag response. It is calculated using the expression n sl t × 100, where n is the number of tags, and sl t denotes the total number of consumed slots. In the case of a collision tags will retransmit after a random delay.
Tags transmit their ID in synchronous time slots.
In case of a collision, tags retransmit after a random delay.
Each tag responds only once per frame.
Tags transmit once per frame. The reader uses a tag estimation function to vary the frame size.
The reader dynamically adjust critical parameter (Q) based on the type of replies from tags.
Disadvantages
In a dense tag environment the the number of collision increases significantly.
In a dense tag environment the the number of collision increases significantly. The reader requires synchronization with tags.
It uses a fixed frame size and does not change the size during the identification process.
It cannot move into the next frame at any time based on the situation of collision without finishing the current frame.
This protocol may encounter some problems (lower throughput) on adjusting Q, especially when the frame size is larger than the number of tags.
Tree-Based Protocols
One of the main features of tree-based protocols is that they are deterministic, since ideally, they are designed to identify the whole set of tags in the interrogation area [36][37][38][39][40]. These protocols have simple design tags and work very well with a uniform set of tags.
Tree-based protocols usually work with a muting capability since they need the identified tags to remain quiet after their identification. These protocols usually work using queries, which are broadcast commands transmitted by a reader to require the tags to respond. If a tag's ID does not match the query, the reader command is rejected.
First, the most popular tree-based protocols are presented here. Then, a selected group of protocols will be presented with the common feature of the use of Manchester coding. The first group includes: Query Tree (QT), Query Window Tree (QwT), and Smart Trend Transversal (STT). Another group consists of tree-based protocols that use Manchester coding: Binary Search (BS), Collision Tree (CT), Optimal Query Tracking Tree (OQTT), and Collision Window Tree (CwT).
Query Tree Protocol
The query tree protocol (QT) is one of the most representative memoryless protocols, in which the reader must provide the tags with a query and the matching tags must respond with their full ID [41]. Tag response depends directly on the current query, ignoring the prior communication history. QT tags involve only simple hardware requirements because they only compare the reader query with their own ID and respond if it coincides. The identification process consists of more rounds in which the reader sends a query, and tags whose ID prefix match the current query respond with their whole ID binary value. In the case of a collision, the reader forms two new queries by appending q with a binary 0 or 1. New queries will be placed in a Last Input First Output stack (LIFO). If there is no response to a query, the reader knows that there is no tag with the required prefix, and the query is rejected. This kind of slot is called idle. If just one tag responds to the reader query, that tag will be identified. By extending the query prefixes until only one tag's ID matches, the algorithm can identify the rest of the tags. The identification procedure is completed when the LIFO stack is empty. Figure 12 shows the QT protocol being used to read 6 tags (Tag A-Tag F). Each tag uses an ID length of k = 6 bits. Initially, the LIFO stack is empty, and the reader begins with a null string. After a collision occurs, the reader pushes queries 0 and 1 into LIFO stack. During the second round, the reader pops from the stack and transmits query 0. In the example in Figure 12, tags 000100 and 001010 match the required prefix, which causes both to transmit and collide. The reader is unable to understand the messages from the tags. The reader then pushes into the stack queries 01 and 00. In the next round, the reader transmits query 00. Again, both protocols respond with their ID and a new collision occurs. In the stack, the following new queries are added: 001 and 000. The reader transmits query 000 and only one tag responds (000100). This tag is successfully identified and will not answer any of the following reader requests. The reader then transmits query 001 in slot 4, which matches tag 001010. In the next round, the reader pops and transmits query 01. For this query, there will be no response since no tags contain that prefix. In round 7 the reader transmits query 1 and the tag from the right side of the tree responds. Four tags will receive this query and a new collision occurs. The reader experiences a collision, since tags 100011, 101110, 110110 and 111001 responded to the query 1. As a result, queries 11 and 10 are pushed onto the stack. The identification process is repeated until round 13, in which the reader transmits the last query (111) from the stack. Overall, the reader uses 13 rounds to read 6 tags.
Smart Trend Traversal Protocol
The Smart Trend Traversal protocol (STT) is a deterministic and memoryless protocol that was created with the aim of reducing the number of collisions in the QT protocol [46]. This protocol has the ability to dynamically issue queries according to an online, learned tag density and distribution. It proposes a combination of the QT protocol and the shortcutting method in order to skip a query that results in a collision. When the protocol detects the potential possibility of a collision, it will avoid it and move to the bottom level of the binary query tree. STT provides trend recognition.
The reader keeps track of the tag density and distribution in order to issue subsequent queries, and consequently, minimizes the number of empty slots and collision slots.
In this protocol, it is not necessary to have any prior knowledge of the network, and it outperforms the existing protocols. The ideal number of queries can be the total number of single nodes. The ideal queries group, referred to as the query traversal path (QTP), is denoted by Q = q 1 , q 2 , q 3 , . . . , q n , where q n is the last query used in the identification process [47].
It is difficult to achieve, but is desirable to get close to its value. The reader can calculate the subsequent queries depending on the tag response, which can be classified into three types: • A collision occurs when the QTP is at too high a certain level and should be moved down by adding a longer prefix to the query. Consequently, the reader appends t bits of 0's to the last query, where t = s + n col − 1. Let s denote the minimum increase, and n col be the number of consecutive colliding slots.
•
An idle slot occurs when no tag responds to a reader query. QTP needs to traverse up just one level, which can lead to a new collision. This rule will be applied only to the right side. If the empty response comes from the left side of the tree, QTP must move horizontally to the right. The reader will decrease the query length by m bits, where m = s + n emp − 1 and n emp is the number of consecutive idle slots. • Upon a successful response, a single node is visited, indicating that the tag has been successfully identified by the reader. Then QTP moves to the symmetric node if the query finishes with a 0, but it returns one level if the query finishes with a 1.
The identification process of the STT protocol, which was explained above, is depicted in Figure 13 with 4 tags.
In conclusion, STT significantly reduces the number of collisions, the identification time, and the energy consumption as compared to the existing Aloha-based and tree-based protocols.
Window Based Protocols
In the majority of tree-based protocols, tags respond with their full ID or with the bits from the last query, when the query sent by the reader matches the tag ID prefix. Figure 14 shows an example of a communication slot between the reader and the tag. To reduce the number of bits transmitted by the tag, a window method has been proposed [48][49][50]. In the identification process, many slots ultimately collide, resulting in a huge waste of bits. Protocols using the window method reduce the number of bits transmitted by the tags. The window is defined as a bit-string of length ws bits transmitted by a tag in a slot. This bit-string is computed on the reader side, respecting the condition 0 < ws < k. It is shown in Figure 15. Most tree-based protocols use a fixed tag response during the identification process, but some use different operational process methods with a dynamic response that is based on window synchronization.
Query Window Tree Protocol
The Query window Tree protocol (QwT) is a memoryless tree-based protocol that applies a dynamic bit window to QT [48,50]. Tags respond directly depending on the current query. QwT tags compare their ID value with the query received and transmit a certain number of bits, managed by the reader. This reduces the complexity of passive tags, their energy consumed, and the identification time. A reader and tag flow chart for QwT are shown in Figure 16a,b. When tags appear in the interrogation area, the reader will broadcast to them by transmitting a query length of L bits. Tags will respond if their ID prefix matches the query sent by the reader, but with the previously specified number of bits. One of the main features is that the total number of collisions is decreased by transforming potential collisions into partial successful slots. This is a new type of slot, called go-on slots. The previously explained window methodology is implemented in the QwT protocol. The window allows tags to transmit only the bit-string instead of their full ID. If tags match a reader query, they will synchronously transmit the next adjacent ws bits of the ID. This protocol uses cyclic redundancy check (CRC) in order to differentiate between the types of tag responses. Accordingly, the slot types that can occur in the QwT protocol can be classified into 4 groups:
•
Collision slot. When the reader cannot differentiate the answer, they will create two new queries by appending '0' and '1' to the former query [q 1 ,q 2 ...q L ]. The window size ws, will remain unchanged, with the value used in the previous query. • Idle slot. When there is no response, the reader will discard the query and retain the same ws as that of the last command.
•
Go-on slot. This occurs when at least one tag responds with a window and the reader is able to understand it. If L + ws < k is not true, the reader will transmit a new query created from the former query and received window. During this query, the reader will append an updated ws value. • Success slot. This is a type of go-on slot where the reader successfully receives the last part of the tag ID and L + ws = k. Then, the reader can save the tag, calculate the new ws, and continue with the identification process.
Using the QwT protocol, the reader computes ws using the expression (1), where β is an adjustable parameter. This heuristic function is used to provide dynamism to the value of ws. It can only be applied to the go-on and Success slots, since in a Collision or Idle slot, ws will be held unchanged. The proposed protocol maintains the memoryless feature of QT in that it is an applied bit window procedure. It provides a decrease in the number of tag-transmitted bits, but increases the number of slots and reader-transmitted bits. Altogether, this tree-based protocol achieves significant energy savings and a reduction in identification time.
• A modification of the QwT is presented in [51] called Standardized Query window Tree protocol (SQwT). This protocol aims to reduce the number of bits that define ws by using a standardized value of 3 bits and approximating it to the nearest power of 2, using s = log 2 ws instead. Tags calculate the number of bits to respond by using ws = 2 s . By using only 3 bits, SQwT can cover window size values from 1 to 128 bits.
•
Another modification of QwT is presented in [52] and is called Flexible Query window Tree protocol (FQwT). This modification takes advantage of the window to perform an estimation of the tag ID distribution in the interrogation area and improve the identification time of the protocol. The functioning of the protocol is divided into two phases: the estimation of the distribution, and the identification process. During the former phase, the reader estimates the tag ID distribution until the first tag is identified (see Figure 17).
After the first tag is identified, the reader begins the identification phase, similar to that of the SQwT. This phase, however uses a different heuristic function when a go-on slot occurs, taking advantage of the c g parameter (see Equation (2) related to the type of tag ID distribution).
Equation (2) is adjusted with a value of parameter β, preselected to decrease the energy consumed by the proposed protocol. Table 2 compares some of the standard tree-based protocols. Defines the generation of the new queries according the several predefined rules. These rules take into account the number of consecutive collisions and slots.
The reader transmits the number of bits tags must respond in addition to the query. This number is calculated using an heuristic function at the reader side.
Disadvantages
It produces a high number of collisions, particularly in the beginning of the identification procedure.
It transmits the full tag ID with each response, therefore a high number of bits are wasted with every collision.
The reader command needs a high number of bits to represent the size of the tags' responses.
System cost Very low Expensive Medium
Complexity Low High Medium
Manchester Coding
Some tree-based protocols work with Manchester coding, which can be used to locate bits that have collided [2,5,53]. The use of Manchester coding to trace the collision to an individual bit is called bit-tracking in the literature [54]. In Manchester coding, the value of a bit is defined by the change in the voltage level: a negative or positive transition. A logical 0 is coded by a positive transition; a logical 1 is coded by a negative transition. In the case in which a minimum of two tags simultaneously transmit bits with different values ('0' and '1'), the positive and negative transitions of the received bits violate the coding rules, and a collision can be tracked.
As shown in Figure 18, Tag 1 is 1100 and Tag 2 is 1010. Both tags synchronously transmit data. The reader can understand the first bit, but the second and the third bit cause a collision. The reader detects a violation of the Manchester codification on those bits, and this is interpreted as a collision located at bits 2 and 3.
Binary Search Protocol
The procedure in the Binary Search protocol (BS) algorithm [2,55] involves transmitting a serial number from the reader to all the tags in the interrogation area. Only tags which have an equal or lower ID value than the received serial number will respond to the request. Then the reader checks the tags' responses bit by bit using Manchester coding and if a collision is detected, the reader divides the tags into subsets based on the collided bits. Table 3 shows an example of BS being used to read four tags (Tag A to Tag D). The reader begins by interrogating tags with the maximum ID value 111. Tags with a value of less than 111 will respond to the query. Their answer results in collision XXX, where all three bits have experienced a collision. In the next slot, the reader transmits a new query by replacing the most significant collided bit (MSB) with a 0. The reader transmits a new query, 011, in the next slot, and all tags compare their ID with the received value. Communication in this slot again results in a collision (01X). In the second slot, the reader replaces the third bit of the command with a 0 and transmits the next query, 010. In the new interrogation round (slot 3) only Tag A has a value equal to or lower than 010, and therefore it is successfully identified. After this slot, the reader restarts the query value with the initial value 111 and transmits it. This procedure is repeated until all of the tags are identified. This protocol has two additional versions: Enhanced BS protocol (EBSA) and Dynamic BS protocol (DBSA) [56]. The main difference from EBSA is that it does not restart the reading procedure after a tag is identified, as in the basic version of BS. To reduce bit consumption, in the initial slot, the reader transmits only '1' instead of all '1's. In the DBSA version, the reader uses the knowledge from the last slot and reduces the number of transmitting bits. For example, if the reader has received 01X, it will request the tags to transmit only the last bit, since the initial prefix has already been identified.
Collision Tree Protocol
The Collision Tree protocol (CT) is an improvement of QT which uses bit-tracking technology in order to find which bits have collided as well as where they are [54,57]. The reader, using the bit-tracking technology, can trace a collision to an individual bit and get the correct bits successfully. This feature works using Manchester coding, which can locate the conflicting bits based on voltage transitions.
The basic features of this protocol is that it decreases collision slots and eliminates idle slots. This contributes to improved results in terms of latency and the number of bits transmitted.The advantage of this protocol compared to the QT protocol is that CT has no idle slots and reduces collision slots. Figure 19 reveals how this protocol works in an environment with 4 tags. At the beginning of the identification process, the reader generates two queries '1' and '0' into a LIFO stack. Then, the reader pops query '0' from the stack and transmits it to the tags. In this case, one tag (010110) matches the query and responds with its ID, and the tag is identified. Then, the reader sends a new query from stack '1' and a collision occurs. Through bit-tracking, the reader can find the colliding bits and thereby resolve potential collisions. The reader pushes two new queries '11' and '10', and firstly transmits '10'. The second tag is identified (101010). On the next transmission, a collision once again occurs. The reader can trace the collision to the fourth bit. Two new queries are made: '1111' and '1110'. These are the last queries in the interrogation round because both tags (111011, 111101) are identified. From this example it may be noted that there are no idle slots and the number of collision slots and latency are reduced, which is the basic aim of the CT protocol.
In conclusion, CT is a stable and efficient anti-collision protocol for RFID tag identification. The performance of CT is very dependent on the total number of tags in the interrogation area.
Optimal Query Tracking Tree Protocol
The Optimal Query Tracking Tree protocol (OQTT) divides all of the tags in the interrogation area into small tag sets in order to reduce the number of collisions at the beginning of the identification process [58]. This protocol uses three main approaches: bit estimation, an optimal partition, and a query tracking tree.
With bit estimation, the reader, using bit tracking technology, can estimate the number of tags in the interrogation area with a small deviation. This phase detects the status of the bits to perform the estimation. The reader broadcasts the parameter l, which denotes the default value of the tag ID length. After receiving a command, all tags must randomly choose a value k between 0 and l−1.
To simplify the procedure, the tags only respond with a bit string of length b, instead of a bit string of length l. All tags generate a b-string of all "0" and set the bit k mod b to "1". Accordingly, the reader can compute the number of selected bits (NSB), and the number of non-selected bits (NNB) from the tags' responses. The probabilities of bits being selected or of being non-selected are calculated from the expressions presented in (2). Finally, the optimal estimation of the number of tags is denoted by ñ and is calculated from (3).
The next approach is an optimal partition which determines the number of initial sets. The reader divides the tags into different sets with initial queries.
The query tracking tree is the last procedure in OQTT, and splits the set of collided tags into two subsets using the first collided bit of the tags' responses. This procedure is followed until no more queries remain in the stack.
The queries in the optimal partition are calculated using the equations c 1 + c 2 = m and 2c 1 + c 2 = 2l, where c 1 denotes the number of (l−1)-bit queries and c 2 the number of l-bit queries. The number of bits in each query is computed from the equations 2l − 1 < m ≤ 2l. Table 4 shows an example of OQTT. In the interrogation area there are 5 tags, whose IDs are as follows: '1010', '1100', '0100', '0010' and '0111'. At the beginning of the frame, the reader estimates the number of tags, which here is n = 5. By using this estimate for the number of tags, the initial number of sets is calculated with the formula m = 0.595824 × n , which yields 3. When m = 3, there are two 2-bit queries and one 1-bit query. The reader generates the queries 00, 01 and 1 and pushes them into the stack. In the presented example, c 2 = 2 (queries 00 and 01) and c 1 = 1 (query 1). As presented in Figure 11, when the reader pops a query, the tags answer with the value k − q, where k and q are the lengths of the ID and the query, respectively. When a collision occurs, this protocol splits the collided query according to the first collided bit. This is the case with slots 3, 4, 5, and 6. However, OQTT may incorrectly estimate the tag number. The estimation error however is negligible, only producing an imperceptible difference in the final result of the number of queries. According to the literature, e.g., [58], this protocol provides an efficiency of approximately 0.614 and is one of the most efficient anti-collision protocols for tag identification. Although the slot efficiency obtained by OQTT is very high, the preprocessing increases the energy consumption of the protocol, especially in dense tag environments [58].
•
A modification of the OQTT is presented in [59] called Optimal Binary Tracking Tree (OBTT). This modification implements the estimator of the OQTT with a simple Binary Tree (BT) protocol [60]. The estimator establishes the inital upper bound for tags' counters. This produces a separation of the existing tags in groups avoiding excessive responses in the beginning of the interrogation procedure.
Collision Window Tree Protocol
The Collision window Tree protocol (CwT) is the second proposed window-based protocol that applies a dynamic window to CT [49,50]. This protocol adopts two techniques: bit tracking and bit windowing. The bit tracking uses Manchester coding in order to identify the colliding bit in the tags' responses. This technique avoids using the CRC, which QwT used in order to identify the type of slot. This protocol does not remove idle slots as CT does, but instead decreases the total amount of bits transmitted by all the tags.
The reader interrogates tags by transmitting a query [q 1 ...q L ] of length L, attached to the ws of length [log 2 ws] + 1 bits. The bit-string ws informs the tags of the number of bits that they must send in their reply. The variable ws is computed in every slot and is transmitted together with the query. Only matching tags transmit ws to the last query bit received, [t L + 1.
When the reader transmits a query, three possible slot statuses can take place after a tag's response: A go-on slot occurs when at least one tag responds and the expression L + ws < k is met. Then, the reader creates a new query based on the former one and the received window from the last slot. The next ws is calculated using the heuristic function in Equation (1). • A success slot occurs when the reader checks the expression L + ws = k and if it matches, the tag is successfully identified.
A CwT flow chart is presented in Figure 20. The reader transmits the first query (0) with ws = 1. The matching tags respond with ws bits and the reader looks for a colliding bit. After slot identification, the reader creates a new query and calculates the value of ws, depending on the type of slot.
The CwT provides a significant decrease in the number of tag-transmitted bits, but this benefit comes with a certain decrease in the number of slots and the reader-transmitted bits. This protocol achieves important energy savings due to the reduction in the time required by the tag transmission process [50]. Table 5 compares some of the tree-based protocols. Table 5. A comparison of tree-based protocols using bit-tracking.
Protocol feature
The reader performs an initial estimation to split the tags in groups. These groups are identified using the CT protocol.
The reader transmits a serial number. Only tags having an equal or lower ID value than the received serial number will respond on request.
It is an improvement of QT which uses Bit tracking technology in order to find which bits collided as well as where they are.
The reader is able to identify the collided bits from the tags' responses, and uses them to calculate the new tag response size.
Disadvantages
Very complex protocol. Hard to be physically implemented.
The reader restarts the reading process after a tag is identified.
It wastes a high number of tag bits on every collision, increasing the energy consumed by the reader during the process.
Hybrid Protocols
Hybrid protocols combine the advantages of tree-based and Aloha-based protocols to avoid their problems and provide better features in tag identification [61][62][63][64]. Most of them first implement a tree-based procedure and tag estimation procedure in order to predict the number of tags. Therefore, the combined Aloha-based and tree-based protocol procedures are known for their high complexity and hardware demands. This kind of protocol can significantly increase performance as compared to the previous ones.
Recent proposals include the Tree Slotted Aloha (TSA) and Binary Tree Slotted Aloha (BTSA). TSA uses a tree structure, and the tag's responses are organized in slots, as in FSA. In the BTSA protocol, tags randomly choose a slot after the reader query.
Tree Slotted Aloha
Tree Slotted Aloha (TSA) is a probabilistic protocol created to reduce the number of collisions occurring in FSA [62]. When more tags collide in a slot, FSA attempts to solve this problem in the next frame. However, in the new approach, if more tags collide in a frame, only those tags that are involved in that collision are queried in the next slot.
TSA uses l o -estimation for the initial frame size. This protocol provides very good efficiency, despite the fact that this number can be far from the actual number of tags.
The initial query consists of a request for data by specifying the frame size l i . Then, all tags in the interrogation area will generate a random number in the range [0, l i ] and transmit its ID in that randomly selected slot.
The protocol is organized in a tree structure. The first node in a tree is the first interrogation round. The reader sets the initial frame with the following data: l 0 , N i represents the number of transmitting tags in slot i, where i ≤ l 0 , N i ≥ 0, and ∑ i N i ≥ n must hold. IfN i ≥ 2, there is a collision in slot i.
At the end of each interrogation round, if the reader detects a collision, it begins a new frame from each slot where the collision was detected. This is accomplished by adding new nodes to the tree: every new node is a son-frame of the collided slot. In each round, the tags store the generated random number from the previous round and increase their tree level counter by 1 so they will know when they should transmit. Every time the reader detects a collision, it creates a new node in the tree and a new round involving only the tags that have collided in that slot. This procedure is shown for the example in Figure 21. TSA is a modified version of the FSA protocol, created to reduce the number of collisions. This protocol behaves better than FSA. TSA achieves an efficiency of between 37% and 41% [62].
Binary Tree Slotted Aloha
The reader in the Binary Tree Slotted Aloha (BTSA) uses a dynamic frame length adjustment and BTSA algorithm [65]. Each tag from the interrogation area randomly chooses a slot and transmit its ID. If the reader successfully identifies a tag it will not be activated in the subsequent slots. When a collision occurs, the collided tags are resolved by binary tree splitting, while the rest of the tags will wait until that process is successfully completed.
The collided tags are continually split into two sets until each set has only one tag. This operation is performed by the BT. The initial frame length is L = 2 Q and the highest efficiency is achieved when the initial frame size is close to the number of tags. Since BTSA has no estimation of the tag set size, the reader cannot set the initial frame size according to the number of tags. Some protocols are presented in order to achieve higher efficiency in a wide range of number of tags. An example of BTSA is shown in Figure 22.
Dynamic Binary Tree Slotted Aloha
Dynamic Binary Tree Slotted Aloha (dynamic BTSA) involves a dynamic frame adjustment and the basic BTSA algorithm [65]. The advantage of this protocol is that the reader can adjust its frame size by judging only the first slot type in the identification process. Figure 23 shows the dynamic BTSA algorithm, where the initial frame length is L = 2 Q , and the initial Q 0 = 4.0. The procedure is very similar to the Q protocol. Firstly, the reader transmits a QueryAdjust command with frame length to all of the tags in the interrogation area. Subsequently, each tag randomly chooses a random number between 0 and L−1. Tags whose Counter value is 0 transmit their ID. Then the reader transmits a new request with a new L and will be capable of receiving responses in the first slot of the following frame. If the first slot is idle, the reader decreases the value of Q by 1 (Q = Q−1) and the reader creates a new frame based on the updated Q. If the reader successfully identifies a tag in the first slot, Q will not be changed and the reader will move to the BTSA algorithm [65]. In BTSA, the reader transmits the Query command in a frame. The reader has a slot counter (SC) that is set to 0 at the beginning of the frame, and is increased by 1 at the end of each slot. When the frame length is equal to SC, the frame finishes. When the reader receives an ID in a slot, it knows the type of slot and will inform other tags by transmitting its feedback. If the reader detects a collision in a slot, it will resolve the collided tags by BinTree splitting [66]. In order for the reader to know when the binary tree has finished, it uses the variable B The initial value of B is set to 2. In the case of a collision, B is increased by 1, and if there is no collision, B is decreased by 1. Only if B = 0 does the reader know that the binary tree has finished.
Dynamic BTSA reduces the number of collisions and improves identification efficiency [65].
Adaptive Binary Tree Slotted Aloha
Adaptive Binary Tree Slotted Aloha (Adaptive BTSA) offers an improvement to the Q protocol [65]. This protocol adjusts the frame size based on the tags' responses in a current slot. Adaptive BTSA first uses features from the Q protocol. If there are numerous collisions in a frame, the reader ends the frame earlier and transmits a new command with a new frame length. If there are excessive idle slots, the reader once again ends the frame earlier and sends a new command with a smaller frame length.
The reader uses the parameters B and Q f p in order to calculate the frame length. The initial frame length is to L = 2 Q and Q = 4. The Q algorithm can adjust the frame length by adjusting Q. The value of Q. is the rounded value of Q f p , which is a floating representation of Q. In the following process, the reader dynamically adjusts each slot using the presented values c [34]. In the first slot, if a collision occurs, the reader will calculate Q f p by increasing it by c. In the case of an idle slot, the reader decreases Q f p by c. When the reader identifies a tag, it will not change Q f p . The flow chart of this protocol is shown in Figure 24.
The function framesize(Q) shown in Figure 24 denotes that a new frame has started and its length is 2 Q . Adaptive BTSA combines the Q algorithm and the BinTree strategy. The main difference between Adaptive BTSA and the Q protocol is that when a collision occurs in the slot, the collided tags will be resolved by BinTree.
Discussion
Anti-collision protocols are a critical part of any RFID system. This section offers a critical analysis of the different protocols presented in the previous sections. The tag collision problem results in the wasting of bandwidth and energy, and an increased latency. Thus, an optimized anti-collision protocol is essential for a competitive RFID system.
The breadth of the literature reveals that there has been a great amount of research conducted in this area. There are two main types of anti-collision protocols: deterministic and probabilistic. In probabilistic protocols, the tags transmit their own ID in randomly selected slots in a frame in order to reduce the possibility of a collision. The tag answers are distributed into the slots and all of them have a chance of being identified. These type of protocols are highly adaptable to the appearances and disappearances of tags from the interrogation area. Deterministic protocols, on the other hand, are ideally designed to identify the whole set of tags in the interrogation area during each cycle. These protocols usually have a simple tag design and can work very well with uniform sets of tags. However, these protocols do not admit unexpected appearances and disappearances as easily as Aloha-based protocols. Tree-based protocols must restart their reading process if a new tag appears in a reader's interrogation area while the tags are being read.
And finally, these types can be combined to form hybrid protocols, which provide very competitive protocols. Hybrid protocols have been created in order to avoid the problems of the Aloha and tree-based protocols, but this comes at the expense of complex reader and tag designs. Table 6 shows observations regarding Aloha, tree-based and hybrid protocols.
From these explanations of the protocols, it cannot be concluded that certain protocol types stand out from the rest. However it should be noted that the newest protocols have become more sophisticated and attain better results in simulations. This also contrasts with the ability to implement these solutions in real hardware. RFID systems are very constrained systems, and their hardware needs to be very simple in order to comply with tags' needs. That is why many of these solutions have yet to be tested under real hardware conditions. Table 6. A comparison of Aloha, Tree-based and Hybrid protocols.
Criterion
Aloha Protocols Tree-Based Protocols Hybrid Protocols
Protocol feature
They use random multi-access means to identify tags. In the case of collision, tags will be asked to send data later with a random time relay.
They identify the total number of tags in the interrogation zone. The reader controls every step of the protocol, using commands or queries to split colliding tags into subsets, and further repeatedly split those subsets until identifying all of the tags.
Tree-based protocols. They use two methods. The first uses randomized divisions in tree-based algorithms, and another uses tree strategies after a collision in Aloha algorithms.
Usage
Aloha protocols are commonly used in LF, HF and UHF RFID (18000-6C) systems.
Not implemented on any standard
Method Probabilistic Deterministic Mixture (Aloha and Tree-based)
Conclusions
The term IoT, as established in RFID and the supply chain, involves the global information service architecture for RFID tags, that is, networked services that discuss things, rather than services. RFID is a key opportunity for the IoT due to its cost-effectiveness, high readability rates, automatic identification and, importantly, its energy efficiency benefits.
This paper presents some of the main RFID procedures and proposes some of the most up-to-date anti-collision protocols. In the literature, these may be classified into Aloha-based, tree-based, and hybrid protocols. The breadth of the literature reveals that there has been considerable research carried out in this field. However, further research needs to be conducted in order to ultimately implement all these solutions. RFID systems have become more and more widely present. Given that the number of tags has increased, these system are faced with more important issues, and therefore, the use of anti-collision protocols will be more omnipresent. | 12,198.2 | 2018-08-01T00:00:00.000 | [
"Computer Science",
"Engineering"
] |
Cognitive Decline, Dementia, Alzheimer’s Disease and Presbycusis: Examination of the Possible Molecular Mechanism
The incidences of presbycusis and dementia are high among geriatric diseases. Presbycusis is the general term applied to age-related hearing loss and can be caused by many risk factors, such as noise exposure, smoking, medication, hypertension, family history, and other factors. Mutation of mitochondrial DNA in hair cells, spiral ganglion cells, and stria vascularis cells of the cochlea is the basic mechanism of presbycusis. Dementia is a clinical syndrome that includes the decline of cognitive and conscious states and is caused by many neurodegenerative diseases, of which Alzheimer’s disease (AD) is the most common. The amyloid cascade hypothesis and tau hypothesis are the two major hypotheses that describe the AD pathogenic mechanism. Recent studies have shown that deposition of Aβ and hyperphosphorylation of the tau protein may cause mitochondrial dysfunction. An increasing number of papers have reported that, on one hand, the auditory system function in AD patients is damaged as their cognitive ability declines and that, on the other hand, hearing loss may be a risk factor for dementia and AD. However, the relationship between presbycusis and AD is still unknown. By reviewing the relevant literature, we found that the SIRT1-PGC1α pathway and LKB1 (or CaMKKβ)-AMPK pathway may play a role in the preservation of cerebral neuron function by taking part in the regulation of mitochondrial function. Then vascular endothelial growth factor signal pathway is activated to promote vascular angiogenesis and maintenance of the blood–brain barrier integrity. Recently, experiments have also shown that their expression levels are altered in both presbycusis and AD mouse models. Therefore, we propose that exploring the specific molecular link between presbycusis and AD may provide new ideas for their prevention and treatment.
INTRODUCTION
The elderly population worldwide is currently approximately 900 million (de Carvalho et al., 2015). There has been a dramatic shift in the distribution of deaths from younger to older ages and from maternal, perinatal, nutritional, and communicable causes to non-communicable disease causes (Mathers and Loncar, 2005). Chronic diseases, of which dementia and presbycusis account for a large part, become more prevalent with age. This trend is exacerbated by lifestyle and behavior changes that predispose individuals to these diseases. According to the World Alzheimer Report, there were an estimated 35.6 million people with AD and other dementias in the year 2010; this number will reach 66 million by 2030 and 115 million by 2050 (Wortmann, 2012). In addition to cognitive decline, AD is also associated with secondary diseases including cardiovascular disease, tumors and sensory system dysfunctions, such as vision and hearing loss (Masters et al., 2015), which impose a heavy burden on patients as well as on society. Therefore, exploring the connections between AD and other diseases is significant for early diagnosis and prevention. As early as 1964, hearing impairment was thought to lead to mental illness caused by isolation (Kay et al., 1964). Subsequently, hearing loss was thought to be independently associated with accelerated cognitive decline and dementia (Lin et al., 2011. Hearing loss is not only the result of auditory organ damage, but may also result from central nervous system dysfunction in auditory information processing. Therefore, an increasing number of researchers have begun to explore whether improving hearing function can improve cognitive disabilities or reduce the risk of dementia later in life. Although Lin et al. (2011) found that treating hearing loss with hearing aids did not significantly decrease the risk of dementia, the risk of incident dementia did increase in participants with a hearing loss of over 25 dB. Dawes et al. (2015) suggested that hearing aids might improve cognitive performance, but this positive effect may not be the result of reducing the adverse effects of hearing loss, such as social segregation or depressed emotion. It remains unclear how presbycusis correlates with cognitive decline, dementia or AD (through social isolation caused by hearing loss, through common neuropathological pathways, or through vascular factors) and whether the relationship between them is unidirectional or bidirectional. As a result, we first review the auditory pathological changes in AD patients and AD mouse models; then, we emphasize how hearing impairment affects the incidence of cognitive decline, dementia and AD on the basis of epidemiologic evidence. Importantly, we summarize the enzymes and proteins that may contribute to the pathogenesis of both presbycusis and AD, including factors such as VEGF, SIRT1-PGC1α, and LKB1 (or CaMKKβ)-AMPK, because AD is the most common type of dementia and the molecular mechanism of AD has been extensively explored. Abnormal expression of these enzymes and proteins may cause AD (Kalaria et al., 1998;Won et al., 2010;Kumar et al., 2013) and dysfunction of cochlear hair cells (Picciotti et al., 2006;Hill et al., 2016;Xue et al., 2016). However, the specific connections between presbycusis and dementia still require further study.
PRESBYCUSIS, DEMENTIA AND ALZHEIMER'S DISEASE Presbycusis
Presbycusis is the general term applied to ARHL. The risk factors for presbycusis include noise exposure, smoking, medication, hypertension, family history and other factors (Mills et al., 2009). According to statistics from the World Health Organization (data from World health Organization [WHO], 2017), a person whose hearing thresholds are over 25 dB in both ears is said to have hearing loss. Hearing loss can be classified as mild, moderate, severe, or profound. Nearly one-third of people over 65 years of age suffer from disabling hearing loss. This disorder is characterized by hearing sensitivity reduction (particularly at high frequencies), reduced understanding of speech in noisy environments, delayed central processing of acoustic information, and impaired localization of sound sources (Tavanai and Mohammadkhani, 2017). Severe hearing loss will affect an individual's psychosocial status and cause social segregation, depression or loss of self-confidence. Although at the early stage, the peripheral auditory dysfunction may be the fundamental pathological change of ARHL, impairment of CAP function becomes increasingly important in late ARHL (Gates and Mills, 2005).
Many molecular theories have been proposed for the development of ARHL. One theory posits that accumulation of ROS in the mitochondria of inner ear hair cells, spiral ganglions and epithelial cells of the stria vascularis can cause further oxidative stress and mtDNA mutations (Yamasoba, 2009;Yamasoba et al., 2013). Markaryan et al. (2009) quantified mtDNA in human cochlear tissue samples, and the deletion rate of mtDNA (4977 bp) was found to reach 32%. However, the deletion rate of mtDNA (4834 bp) was found by Kong et al. (2006) to reach over 90% in the inner ear of rats treated with D-galactose, and this type of mutation can enhance the sensitivity of the inner ear to an aminoglycoside antibiotic.
Dementia and Alzheimer's Disease
Dementia is a chronic and progressive deterioration disease characterized by cognitive dysfunction and abnormal mental behavior. It has become the greatest global challenge for health and social care in the 21st century (Livingston et al., 2017). Two of more common types of dementia are AD, which is characterized as a progressive, unremitting and neurodegenerative disorder (Masters et al., 2015); and vascular dementia (VD), which is mainly caused by hypertension and arteriosclerosis. AD is the main type of dementia that has caught global attention. The population of people with AD and other dementias will reach 66 million by 2030 and 115 million by 2050 (Wortmann, 2012). Dementia is associated with age, and the incidence of dementia increase from 3.9/1000 person years (pyr) at an age of 60-64 to 104.8/1000 pyr at an age of over 90, which indicates that the incidence doubles every 6.3 years. Currently, Aβ, tau protein and ApoE are the three main elements that are thought to contribute to AD. The pathology of dementia, especially AD, which is focused on in the following section, includes: (1) loss of neurons in the temporal lobes and hippocampus; (2) NFTs; (3) SPs that consist of amyloid-β (Aβ); and (4) amyloidopathy of the cerebrovasculature (Hyman, 1998). NFTs and SPs are characteristic pathological changes. The amyloid cascade hypothesis and tau hypothesis are the two major hypotheses that describe the AD pathogenic mechanism. The average duration of illness is 8-10 years, but the clinically symptomatic phases are preceded by preclinical and prodromal stages that typically extend over two decades (Masters et al., 2015). Therefore, it is of great clinical significance to diagnose and prevent AD at an early stage.
PATHOLOGICAL CHANGES OF THE AUDITORY SYSTEM IN PATIENTS WITH AD
Central auditory processing dysfunction is highly evident in persons with Alzheimer's disease (Idrizbegovic et al., 2011), and pathological changes have also been found in the auditory system, as described in the following studies. The auditory nervous pathway originates from afferent neurons called spiral ganglion cells in cochlea. Spiral ganglion cells are a type of bipolar neurons located in the Rosenthal canals of the bony modiolus. Next, the axons of spiral ganglion cells project to the cochlear nucleus complex. Then, most of the axons from the cochlear nucleus complex cross the midline and ascend in the contralateral lateral lemniscus, terminating in the inferior colliculus and medial geniculate body. Finally, all ascending neurons form an auditory radiation and terminate in the auditory center of the transverse temporal gyrus (Lonsbury-Martin et al., 2009). Neurons in different parts of the medial geniculate body are represented with the low best frequencies arranged laterally and high best frequencies arranged medially (Aitkin and Webster, 1971). Volume and quality loss of the cerebrum (meaning the sulcus is wider and deeper) as well as atrophy of the gyri have been observed in the cerebrum of AD patients using imaging methods. Consequently, we hypothesize that patients with AD might have hearing loss at both low and high frequencies if the neurons in the medial geniculate body are widely degenerated. Therefore, every link in the auditory pathway may display pathological changes associated with AD, thus leading to hearing impairment (Figure 1).
Cochlear Pathology
The production and conduction of auditory signals in the cochlea mainly involve hair cells and spiral neurons. Wang and Wu (2015) reported that a significant loss of SGNs, rather than hair cells, could be found in the cochlea of 9-and 12month-old 3xTg-AD model mice. However, O'Leary et al. (2017) found that in 5xFAD model mice (a kind of AD model mice), outer and inner hair cells also showed significantly greater losses at the apical and basal ends of the basilar membrane than wild-type mice at 15-16 months of age. Omata et al. (2016) produced new transgenic mouse models [Tg(Math1E-Aβ42Arc)1Lt/Tg(Math1E-Aβ42)1Lt], which overexpress Aβ or Aβ-related proteins in cochlear hair cells, and Tg (MathE-MAPT) 1Lt, which express human tau (2N4R) in cochlear hair cells. These transgenic mice showed auditory dysfunction, especially in high-frequency sound perception, and in these mice, expression of Aβ and tau was correlated with hearing impairment as well as hair cell loss. Consequently, it is speculated that abnormal deposition of Aβ and overexpression of tau protein in cochlea hair cells may have a synergistic effect on hearing impairment.
Pathology of the Medial Geniculate Body
The ventral nucleus of the medial geniculate body is one of the most important relay stations in the ascending auditory pathway, and it receives fibers from neurons in the central nucleus of the inferior colliculus. Sinha et al. (1993) found that SPs and NFTs were extensively distributed throughout not only the ventral nucleus of the medial geniculate body but also the central nucleus of the inferior colliculus in AD patients at autopsy. Similarly, Rüb et al. (2016) found conspicuous AD-related cytoskeletal pathology in the inferior colliculus, superior olive and dorsal cochlear nucleus. Parvizi et al. (2001) also found β-amyloid and hyperphosphorylated epitopes of the tau protein in the inferior and superior colliculus and autonomic, monoaminergic, cholinergic, and classical reticular nucleus.
Auditory Cortex Pathology
The human auditory cortex is mainly located in the superior temporal gyrus (Hackett, 2015). In the early stage of AD, brain atrophy involved the temporal lobe, especially the hippocampus and other brain regions, including the central auditory cortex and its related functional nuclei. In addition to the relay stations in the auditory pathway, such as the medial geniculate body and inferior colliculus, SPs and NFTs have also been observed in the primary auditory cortex and association area of the auditory cortex (Sinha et al., 1993). Recently, expression of VEGF was found to be reduced in the superior temporal, hippocampal, and brainstem regions of AD patients (Provias and Jeynes, 2014). VEGF is an important endothelial growth factor that is responsible for vascular angiogenesis, remodeling, and maintenance of the blood-brain barrier (Sondell et al., 1999). It stimulates axonal outgrowth, thus promoting cell survival and Schwann cell proliferation in the peripheral nervous system (Rosenstein et al., 2010). Moreover, VEGF repairs hair cell damage caused by noise, drugs, or certain diseases, such as otitis media and acoustic neuroma (London and Gurgel, 2014). The results mentioned above imply that reduced expression of VEGF may cause abnormalities in the structure and function of blood vessels and neurons in the auditory cortex of patients with AD, leading to a severe hearing loss. In addition to the pathological changes in auditory pathways, mouse models and patients with AD also showed increased ABR thresholds, and a greater hearing loss was related to higher adjusted relative odds of having dementia (Uhlmann et al., 1989;Goll et al., 2011;O'Leary et al., 2017).
ARHL MAY BE A RISK FACTOR FOR COGNITIVE DECLINE, DEMENTIA OR AD
Defects in sensory systems, including the olfactory, visual or auditory systems, are thought to be highly associated with age-related neurodegenerative diseases (Benarroch, 2010). Impairments in peripheral and central auditory organs have been linked to accelerated cognitive decline (Bernabei et al., 2014;Amieva et al., 2015), incident cognitive impairment (Deal et al., 2017), dementia and AD (Gallacher et al., 2012;Panza et al., 2015;Taljaard et al., 2016). Therefore, audiometric testing may serve as a useful method for evaluating cognitive function, dementia and AD. After a 12-year follow-up study, Lin et al. (2011) noted that hearing loss was independently related to incident dementia after eliminating the influence of age, sex, race, education, hypertension, and other factors in 639 participants aged from 36 to 90 years old, and the attribute risk of dementia related to hearing loss reached 36.4%. In 2013, they further proposed that hearing loss was independently associated with accelerated cognitive decline . In addition to the peripheral auditory function, Gates et al. (2010) proposed that central auditory dysfunction and executive dysfunction might also give rise to neurodegenerative processes. CAP dysfunction in one ear was related to a sixfold increase in the risk of cognitive decline in later life after a 6year follow-up study (Gates et al., 1996). Panza et al. (2015) concluded from many clinical trials that a CAP deficit might be an early sign of cognitive decline, and that CAP testing could be used to evaluate cognitive function in the future. Other researchers believe that hearing defects may play a part in producing mental symptoms, such as social isolation and loneliness, by reducing an individual's contact with the outside world due to communication impairments caused by ARHL (Bennett et al., 2006). Evidence from epidemiological studies suggests that hearing loss is a modifiable factor, and with the help of appropriate and timely treatment, cognitive decline can be decelerated and daily activities can be facilitated (Behrman et al., 2014). Cognitive performance was found to be improved in auditory rehabilitation with the use of cochlear implants or hearing aids, which suggests that interventions that aim to restore hearing may be an effective way to alleviate cognitive disorders in late life (Mulrow et al., 1990;Castiglione et al., 2016). Individuals with greater hearing loss might receive the most cognitive benefit from hearing aids (Meister et al., 2015). Although many researches have shown that hearing loss may increase the risk of cognitive decline, as mentioned above, there is still much debate regarding whether there is a link between hearing loss and dementia. Gennis et al. (1991) found that hearing loss had no significant relation with cognitive function after a 5-year follow-up of 224 people aged over 60 years with no serious underlying disease. In addition, some studies have also found no evidence of improvement in behavioral symptoms, functional status, or quality of later life by providing hearing aids to hearing-impaired AD patients Nguyen et al., 2017).
In addition to the above clinical trials, patients with hearing loss have also exhibited significant reductions in the gray matter volume of the cortex related to hearing, attention and emotion, as revealed by morphometry, functional magnetic resonance (fMRI) imaging or EEG studies (Wong et al., 2010;Peelle et al., 2011;Eckert et al., 2012;Cardin, 2016). With fMRI, Peelle et al. (2011) found that a decline of peripheral auditory acuity not only led to a loss of gray matter volume in the primary auditory cortex but also contributed to down regulation of neural activity in the course of language processing. EEG studies have revealed that the greater the degree of hearing loss, the more sluggish the cortical response (Cardin, 2016). These results suggest that changes in the anatomy and function of the brain associated with hearing loss might be a reason for the increased incidence of cognitive decline and dementia. In animal experiments, Yu et al. (2011) found that when C57BL/6J mice developed profound hearing loss by 42-44 weeks of age, their cognitive function also declined, as evaluated by the Morris water maze and that the ultrastructure of the synapses of the CA3 region of the hippocampus also changed, including an increase in the synaptic cleft width and decrease of the thickness of the postsynaptic density. Therefore, we believe that there may be a certain connection between ARHL and dementia. The pathological changes of the hippocampus associated with hearing loss may lay the foundation for the incidence of cognitive decline and dementia. Similarly, Wayne and Johnsrude (2015) also proposed that there may be a common mechanism for both diseases and, to some extent, hearing loss may be an early manifestation of the underlying pathological changes.
ARHL AND AD MAY RESULT FROM MITOCHONDRIAL DYSFUNCTION AND CHANGES IN CERTAIN SIGNAL PATHWAYS
The role of mitochondrial dysfunction in the pathogenesis of dementia is hotly debated. Manczak et al. (2006) found that oligomeric forms of Aβ are mainly aggregated in the mitoplast (the inner membrane plus the matrix) of mitochondria inside neurons of the cortex and hippocampus. In a transgenic mouse model of AD, the hydrogen peroxide levels were greatly increased and cytochrome c oxidase activity was decreased in a transgenic mouse model of AD. These changes were directly associated with the levels of soluble Aβ, which suggest that soluble Aβ may cause more production of hydrogen peroxide and impair mitochondrial metabolism in the development and progression of AD. By contrast, aggregation of ROS may cause mtDNA mutations or deletions, and mitochondrial dysfunction is correlated with the development of ARHL (Yamasoba, 2009;Yamasoba et al., 2013). In recent decades, an increasing number of researchers have begun to pay attention to the function of the mitochondria in an effort to explore whether there is a common pathway in mitochondrial oxidative metabolism that can cause both dementia and hearing loss. If there is a common pathway between them, it may provide new ideas for preventing and treating AD as well as hearing loss.
ROS/VEGF Pathway
Reactive oxygen species, such as O 2− and H 2 O 2 , can activate VEGF, which is important for angiogenesis and neuron protection. VEGF transmits signals to endothelial cells mainly through two tyrosine kinase receptors, VEGFR-1 (Flt-1) and VEGFR-2 (Flk-1) (Rosenstein et al., 2010). Research has shown that flt-1(-/-) embryos have defective sprouting from the dorsal aorta (Kearney et al., 2004). Exogenous H 2 O 2 administered to human umbilical vein endothelial cells to simulate ROS production in cells increases mitochondrial (mtROS) through Serine 36 phosphorylation of p66Shc (Kim et al., 2017). Mitochondria in endothelial cells are also implicated in ROS signaling transactivation of VEGFR2 or VEGF-induced cell migration. Inhibition of mitochondrial respiration with rotenone and oligomycin can attenuate H 2 O 2 -induced tyrosine phosphorylation of VEGFR2 in bovine aortic endothelial cells (Chen et al., 2004). When cells are faced with a hypoxic environment, both VEGF production and VEGFR2 expression are upregulated and the generation of hydrogen peroxide is triggered by mitochondria under hypoxic conditions (Chandel et al., 1998;Colavitti et al., 2002). VEGF may also promote neurogenesis, neuronal patterning, neuroprotection and glial growth by acting on blood vessels. VEGF was found to be upregulated when the central nervous system suffered from injuries, and exogenous application of VEGF facilitated central nervous system angiogenesis (Rosenstein et al., 2010). Wang et al. (2007) found that neurogenesis and neuromigration were enhanced under ischemic condition in VEGF-overexpressing transgenic mice. In addition, the use of VEGF-specific antibodies that blocked its normal function in a stab injury model led to an increase in lesion size and reduction of angiogenic and astroglial activity in the striatum (Krum and Khaibullina, 2003).
As mentioned above, VEGF levels are reduced in the superior temporal, hippocampal, and brainstem regions in AD patients (Provias and Jeynes, 2014). VEGF is also decreased in aged ARHL mice (Picciotti et al., 2004). Therefore, as a downstream signal of mitochondrial regulation, VEGF is likely to be a common factor in the molecular mechanisms of dementia and presbycusis. Since VEGF has the function of neuroprotection, it also exerts many positive effects during the pathogenesis of AD. Expression of VEGF in the brain can increase transiently in the early stage of AD, implicating some compensatory mechanisms to counter pathological changes, such as insufficient vascularity and reduced perfusion (Kalaria et al., 1998). Exogenous VEGF was administered to transgenic mice by Burger et al. (2009) who found that VEGF modulated β-secreatase1 (BACE1) and reduced soluble Aβ1-40 and Aβ1-42. In addition, VEGF was found to bind to amyloid plaques with high affinity, most likely causing a deficiency of available and free VEGF under hypoperfusion conditions and possibly contributing to neurodegeneration and vascular dysfunction in the progression of AD (Yang et al., 2004). The decreased VEGF levels in the brain may affect the function of memory and cognition. The cognitive function of APP/PS1 mice was improved after intraperitoneal injection of VEGF (Wang et al., 2011). Intra-hippocampal administration of a VEGF receptor blocker may negatively affect long-term memory (Pati et al., 2009). Donepezi, an acetylcholinesterase inhibitor, exhibits beneficial effects in Alzheimer's disease through activation of the PI3K/Akt/HIF-1α/VEGF pathway (Kakinuma et al., 2010). By reviewing the relevant literature, we found that VEGF is not only associated with the occurrence of dementia, but may also play an important role in many inner ear diseases. In studies, VEGF has been shown to contribute to the development of vestibular schwannoma and otitis media with effusion, and the use of a VEGF antagonist can alleviate hearing loss in these two diseases (Lim and Birck, 1971;Koutsimpelas et al., 2007;Plotkin et al., 2009;Cheeseman et al., 2011). However, the specific function of VEGF in ARHL requires further investigation. Among the many risk factors of ARHL, exposure to noise and ototoxic drugs, such as kanamycin or cisplatin, has been reported to lead to significant cochlear vascular changes, including an increase in vascular permeability, alterations in cochlear blood flow and vasoconstriction (Hawkins et al., 1972). Acoustic trauma was also found to structurally change blood vessels by disrupting the cochlear blood-barrier (Shi, 2009). In a study exploring noise-induced sensorineural damage, upregulation of VEGF was observed in the stria vascularis, spiral ligament and SGNs of guinea pigs as early as 1 day after noise exposure, and it also preceded the gradual recovery of inner ear function. Therefore, VEGF may also enhance neuron survival during the course of ARHL (Picciotti et al., 2006).
In short, aging often causes reduced ROS clearance and increased ROS activity (Kaur et al., 2016). Accumulation of ROS in cells can cause mtDNA mutations and mitochondrial dysfunction and then result in decreased levels of VEGF. Decreased levels of VEGF may lead to Aβ deposition and hamper the repair of hair cells and spiral ganglions in the inner ear as well as the reconstruction of blood vessels in the brain, resulting in dementia or presbycusis. Further studies are needed to determine whether there is a time order of the occurrence of both diseases and whether hearing loss caused by a decrease of VEGF will further increase the risk of dementia and its progression (Figure 2).
SIRT1/PGC-1α (or FNDC5) Pathway
Amyloid plaques primarily consist of Aβ1-42 and Aβ23-35 (Glenner and Wong, 1984;Gruden et al., 2007). BDNF is a versatile and multifunctional growth factor that is indispensable in a wide range of adaptive processes during human brain development, ranging from regulation of synapse formation and plasticity to neuronal differentiation and better cognitive function (Phillips et al., 1991;Park and Poo, 2013). BDNF precursors and mature BDNF were found to be reduced in the brain of patients in the pre-clinical stages of AD (Peng et al., 2005). Xia et al. (2017) found that both the protein and mRNA levels of BDNF were reduced in APP/PS1 transgenic mice as well as in the hippocampus and cerebral cortex of C57BL/6 mice after injection of Aβ1-42 oligomer. BDNF exerts protective effects in the process of neurodegeneration. In the case of AD, BDNF FIGURE 1 | Auditory signal pathway and possible insults to the auditory nervous system caused by Alzheimer's disease. The auditory nervous pathway originates from spiral ganglion cells in cochlea. Then, the axons of spiral ganglion cells project to the cochlear nucleus complex. Next, most of the axons from the cochlear nucleus complex cross the midline and ascend in the contralateral lateral lemniscus, terminating in the inferior colliculus and medial geniculate body. Another part of axons project to the superior olivary nucleus and then terminate in the inferior colliculus and medial geniculate body. Finally, all ascending neurons form an auditory radiation and terminate in the auditory center of the transverse temporal gyrus. Every link in the auditory pathway may display pathological changes associated with AD, thus leading to hearing impairment.
shifts APP processing toward the α-secretase pathway rather than the β-secretase pathway to repair the neurotoxic effects of Aβ in the brain (Holback et al., 2005). Moreover, BDNF activation is responsible for the tau dephosphorylation, thus inhibiting aggregation of the tau protein in NFTs (Murer et al., 1999;Elliott et al., 2005). A previous study showed that PPARγ coactivator-1 (PGC-1α) and its downstream activated membrane protein, FNDC5, could modulate the BDNF levels together in the PGC-1α (−/−) mouse model (Wrann et al., 2013).
PGC-1α is important in the control of cellular energy metabolic pathways (Finck and Kelly, 2006). It is more highly expressed in cells that are rich in mitochondria and is one of the factors involved in mitochondria synthesis and respiratory gene regulation (Kelly and Scarpulla, 2004). Moreover, PGC-1α also plays an important role in inhibiting neurodegeneration. PGC-1α was reported to protect against the neurodegenerative effects induced by MTPT (1-methyl-4-phenyl-1, 2, 3, 6tetrahydropyridine) and enhance ROS detoxification (St-Pierre et al., 2006). In addition, PGC-1α can promote spinogenesis and synaptogenesis in cultured hippocampal neurons (Cheng et al., 2012). SIRT1 is a NAD+-dependent protein deacetylase located upstream of PGC-1α. It plays a crucial role in maintaining cellular homeostasis by regulating neuron survival and death, glucose metabolism, insulin sensitivity, and mitochondrial synthesis (Lagouge et al., 2006;Guarente, 2013). PCG-1α deacetylation mediated by SIRT1is required for the activation of mitochondrial fatty acid oxidation genes (McCarty et al., 2015).
Reduced levels of PGC-1α and FNDC5 were also found in Neuro-2a cells treated with Aβ oligomers in vitro (Xia et al., 2017). However, when PGC-1α and FNDC5 were overexpressed, the reduction in BDNF caused by Aβ oligomers was reversed, suggesting that the inhibitory action of Aβ on BDNF levels was mediated by a PGC-1α-/FNDC5-dependent pathway. With the help of BDNF treatment, deposition of Aβ in the brain was restrained and the cognitive decline was postponed in APP/PS1transgenic mice. SIRT1 can modulate Aβ metabolism by regulating the processing of amyloid precursor protein (APP) in AD progression (Kumar et al., 2013). Wang et al. (2013) found that treatment with the SIRT1-activating agent resveratrol prevented the generation of Aβ through a SIRT1dependent mechanism. They also showed that both in vitro and in vivo experiments, downregulation of PGC-1α or PGC-1β gene transcription using specific small interfering RNA (siRNA) resulted in augmented expression of β-secreatase1 (BACE1), a protein responsible for cleaving APP. In addition, overexpression of PGC-1α decreased BACE1 protein expression in the hippocampi of transgenic mice (Wang et al., 2013). The dysfunction in the SIRT1-PGC-1α pathway not only linked to AD, but also influenced the function of hair cells and caused hearing loss. Expression of SIRT1, SIRT3, and SIRT5 was reduced in the cochlea of 22-month-old CBA/J mice (Takumida et al., 2016). A previous study revealed that SIRT1 was the direct target gene of miR-29 (Xu et al., 2014) and that expression of miR-29b was upregulated in the ARHL mouse FIGURE 2 | Reactive oxygen species (ROS)/VEGF pathway. Aging often causes reduced ROS clearance and increased ROS activity. Accumulation of ROS in cells can cause mtDNA mutations and mitochondrial dysfunction and then result in decreased levels of VEGF. Decreased levels of VEGF may lead to Aβ deposition and hamper the repair of hair cells and spiral ganglions in the inner ear as well as the reconstruction of blood vessels in the brain, resulting in dementia or presbycusis. Further studies are still needed to determine whether there is a time order of the occurrence of both diseases and whether hearing loss caused by a decrease of VEGF will further increase the risk of dementia and its progression.
model (Zhang et al., 2013). Xue et al. (2016) found an agedependent downregulation in SIRT1 and PGC-1α protein levels, as well as mitochondrial dysfunction in the cochlea of aged mice. In vitro experiments showed that overexpression of miR-29b in HEI-OC1 cells inhibited SIRT-1 and PGC-1α protein levels, causing mitochondrial dysfunction and hair cell apoptosis, whereas in miR-29b knockdown cells, the SIRT1 and PGC-1α protein levels, as well as their mRNA levels, were all significantly upregulated (Xue et al., 2016). Additionally, overexpression of SIRT1 markedly promoted PGC-1α expression in HEI-OC1 cells, inhibited cell apoptosis, and boosted cell proliferation. Consequently, it was hypothesized that miR-29b/SIRT1/PGC-1α signaling most likely plays a role in regulating hair cells apoptosis and in the pathogenesis of ARHL.
From the above research, we propose that expression of miR-29b may be increased in the inner ear with aging, which reduces the expression of its target SIRT1 and affects downstream PGC-1α and FNDC5. Reduced expression of SIRT1-PGC-1α may impair the synthesis and respiratory function of mitochondria, thereby causing neuronal degeneration and cell apoptosis, thus leading to presbycusis. At the same time, reduced expression of SIRT1-PGC-1α in the brain can cause increased production of Aβ and reduced production of BDNF and thus contribute to the incidence of AD. Though no previous study has explored the changes of miR-29b in the brains of AD patients, we believe that SIRT1/PGC-1α (FNDC5) may be a possible drug target to prevent patients with ARHL from developing AD. However, SIRT1 has recently been reported to play a completely adversarial role in the early onset of ARHL in C57BL/6 mice (Han et al., 2016). Whether SIRT1 has a positive or negative effect on ARHL requires further investigation (Figure 3).
LKB1 and CaMKKβ/AMPK Pathway
AMP-activated protein kinase (AMPK) is a major regulator of cellular energy homeostasis and a central player in glucose and lipid metabolism (Daval et al., 2006;Hardie, 2008). First, when the human body is faced with hypoxia, oxidative stress or ischemia, AMPK is activated and phosphorylated by LKB1 complex or CaMKKβ, accompanied by microenvironment changes, such as an elevated AMP/ATP ratio and increased Ca 2+ levels, respectively, thus promoting cellular glucose uptake, fatty acid β-oxidation, glucose transporter 4 synthesis and mitochondrial synthesis (Viollet et al., 2009). Activation of AMPK can lead to more fatty acids entering the mitochondria as well as stimulate the synthesis of ATP (Saha and Ruderman, 2003). It can also promote ATP synthesis by suppressing FIGURE 3 | SIRTl/PGC-1α (FNDC5) pathway. Reduced expression of SIRT1-PGC-1α caused by many factors may impair the synthesis and respiratory function of mitochondria, thereby causing neuronal degeneration and cell apoptosis, thus leading to presbycusis. At the same time, reduced expression of SIRT1-PGC-lα in the brain can cause increased production of Aβ and reduced production of BDNF and thus contribute to the incidence of dementia. Therefore, SIRT1/PGC-lα (FNDC5) may be a possible drug target to prevent patients with ARHL from developing AD.
upregulation of ATPase inhibitory factor 1 (IF1) and stimulating the activity of the oxidative respiratory chain (Vazquez-Martin et al., 2013). Second, AMPK can phosphorylate the PGC-1α protein at threonine-177 and serine-538 to directly regulate mitochondrial synthesis (Jager et al., 2007) or deacetylate PGC-1α via SIRT1activation to modulate mitochondria and lipid utilization genes (Canto et al., 2010). Third, AMPK can improve mitochondrial autophagy due to reactive oxygen damage (Pauly et al., 2012). Therefore, the LKB1 (or CaMKKβ)/AMPK pathway and its association with the SIRT1/PGC-1α pathway are closely related to the function and synthesis of mitochondria.
AMP-activated protein kinase has also been shown to play a key role in regulating neurodegenerative diseases via autophagy of mitochondria (Vingtdeux et al., 2011). Transient glutamate exposure during ischemia or stroke leads to rapid and transient AMPK activation with an increase in Glucose Transporter 3 trafficking, which exerts a neuroprotective effect by increasing the ATP/AMP ratio and decreasing cytosolic Ca 2+ levels (Weisova et al., 2011). In the progression of AD, AMPK activation may downregulate the generation of Aβ by modulating APP processing (Won et al., 2010). At the same time, AMPK is also a physiological tau kinase, and direct stimulation by AICAR (AMPK activator) inhibits tau phosphorylation (Chakraborty et al., 2008;Greco et al., 2009a,b). Moreover, AMPK activation can suppress the mTOR signaling pathway and then enhance cell autophagy and lysosomal degradation of Aβ (Anekonda et al., 2011;Vingtdeux et al., 2011). Sodium hydrosulfide was found to alleviate cell apoptosis in the auditory cortex through the CaMKKβ/AMPK pathways, suggesting that AMPK is a crucial factor in neuron degeneration of the auditory cortex in the central nervous system (Chen et al., 2017). Therefore, AMPK may be a neoteric determinant in neurodegenerative diseases and one of the possible targets for drugs aimed at AD and ARHL. Another upstream kinase of AMPK, LKB1, contributes to maintaining the development and structural stability of cochlear hair cells and stereocilia. Men et al. (2015) found that LKB1 is required for the development and maintenance of hair cell stereociliary bundles and is expressed from the cuticular plate to the nuclei of hair cells. The ABR and DPOAEs thresholds were significantly higher in Atoh1-LKB1-/-mice than those in control mice, although there were no obvious differences between mutant mice and control mice in gross morphology (Men et al., 2015). However, whether expression of LKB1 will change during the aging process and affect the process of ARHL is still unknown.
AMP-activated protein kinase is capable of inhibiting c-Jun N-terminal protein kinase (JNK) activity in neurons (Schulz et al., 2008). However, prolonged elevation of phosphorylated (p)-AMPK can trigger chronic activation of JNK, thereby causing upregulation of the proapoptotic protein Bim (Bcl-2 interacting mediator of cell death) and subsequently leading to apoptosis in neuronal and pancreatic cells, suggesting that cell fate regulation by AMPK is complex (Kefas et al., 2004;Yun et al., 2005;Weisova et al., 2011). Hill et al. (2016) proposed that longterm phosphorylation of AMPK caused by noise exposure could inversely lead to a decreased auditory function together with hair cell and synaptic ribbon loss. After CBA/J mice were exposed to 98 dB of noise, immunolabeling for p-AMPKα showed a stronger change in IHCs and marginal elevation in OHCs, while no obvious changes were detected in pillar cells. In non-exposed control animals, p-AMPKα was weak in hair cells and stronger in pillar cells. As a result, we can speculate that activation of AMPK may exert a dual influence, which may have an opposite effect on cell death and survival.
Therefore, we can hypothesize that decreased activity of AMPK will lead to abnormalities in mitochondrial synthesis and normal functions, such as oxidative respiration, as well FIGURE 4 | Liver kinase B1 (LKB1) and CaMKKβ/AMPK pathway. Decreased activity of AMPK will lead to abnormalities in mitochondrial synthesis and normal functions, such as oxidative respiration, as well as alterations of autophagy, which render cells unable to handle changes in the microenvironment. Decreased activity of AMPK will also cause less degradation of Aβ and more phosphorylation of tau. LKB1 and CaMKKβ are required for the activation of AMPK and are important for the stability of stereocilia in hair cells and protection of auditory cortex neurons, respectively. Although there is currently no evidence showing that LKB1 and CaMKKβ change in the ear or brain due to aging, the environment, diet or other factors, we speculate that if expression of LKB 1 and CaMKKβ is altered by the above factors, they would not only cause hearing loss but also affect downstream phosphorylation of AMPK, leading to AD. However, the activation of AMPK may exert a dual influence, which may have an opposite effect on cell death and survival.
as alterations of autophagy, which render cells unable to handle changes in the microenvironment. Decreased activity of AMPK will also cause less degradation of Aβ and more phosphorylation of tau. LKB1 and CaMKKβ are required for the activation of AMPK and are important for the stability of stereocilia in hair cells and protection of auditory cortex neurons, respectively. Although there is currently no evidence showing that LKB1 and CaMKKβ change in the ear or brain due to aging, the environment, diet or other factors, we speculate that if expression of LKB1 and CaMKKβ is altered by the above factors, they would not only cause hearing loss but also affect downstream phosphorylation of AMPK, leading to AD. In addition, the dual effect of AMPK leads us to suggest a novel concept: AMPK may be phosphorylated in the brains of AD patients for a long time in response to intracellular environmental changes, which may damage the auditory cortex and accelerate hearing loss. Therefore, LKB1 and CaMKKβ/AMPK may also be a new target for drugs to protect patients with ARHL from dementia (Figure 4).
The above-mentioned publications may guide us in studying the relationship between dementia and ARHL. We speculate that expression or activity of SIRT1/PGC-1α and LKB1/AMPK may be changed by factors such as aging, the environment and medication. The expression or activity change may cause dysfunction of mitochondrial synthesis, oxidative respiration, and autophagy, leading to abnormalities of mtROS signal transduction and activation of downstream VEGFR2 transcription, followed by the development of both dementia and presbycusis. These two diseases may also influence the progression of each other as they have common pathological pathways and common targets. Undoubtedly, more research is necessary to confirm the role of these pathways. If we can find a common mechanism and verify that hearing loss is an early manifestation of dementia, it would be meaningful for the early diagnosis and prevention of dementia.
CONCLUSION
Due to the non-renewable characteristic of neurons, there are no completely effective methods for curing dementia or presbycusis. With the increasing aging population, the number of people suffering from AD or other dementias and hearing disability will also increase. The factors that influence dementia and agerelated hearing impairment are highly complex and involve not only age but also the environment, genetics, lifestyle, drugs and other factors. As a result, considering how to delay or prevent the incidence of dementia and how to delay the neurodegeneration in patients with dementia is important. Although pathological changes of the auditory system have been observed in patients or mice with AD, as mentioned in Section "Pathological Changes of the Auditory System in Patients with AD, " and a group of epidemiological studies have revealed that ARHL may be a risk factor for cognitive decline, dementia or AD, as mentioned in Section "ARHL May be a Risk Factor for Cognitive Decline, Dementia or AD, " the specific molecular mechanism linking these diseases is still unknown. It is also still not clear whether the relationship is unidirectional or bidirectional or if they are both the clinical manifestations of aging. Therefore, in vivo experiments are necessary if we are to explore the relationship between AD and ARHL at the molecular biology level. For example, we may make full use of transgenic AD mouse models, such as APP/PS1 mice, 3xTg-AD mice and 5x-FAD mice, to explore whether AD can cause more severe auditory dysfunction compared to non-transgenic mice in the same background. To explore whether ARHL can cause cognitive changes, C57BL6/J mice (a strain that develops late-onset sensorineural hearing impairment) (Noben-Trauth et al., 2003) and CBA/CaJ mice (a strain that has normal hearing throughout life) are good experiment and control groups, respectively. To our expectation, pathological changes have been found in some AD mouse models and cognitive decline has been verified in old C57BL6/J mice compared to CBA/CaJ mice of the same age, according to existing research. Although these are only patterns of manifestations, they support our belief that certain correlations really exist between ARHL and cognitive decline, dementia or AD. The next step is to explore the molecular mechanism, and with the help of a series of molecular biological techniques, we can try to validate whether the expression of proteins such as VEGF, Sirt1, PGC-1α, LKB1, CaMKKβ or AMPK has the same tendency of variation in the peripheral auditory system, central auditory system, temporal lobes and hippocampus. However, we need to overcome obstacles, such as the long time period needed for experiments, since AD and ARHL model mice require at least 6 months or more to exhibit their phenotype. In addition, we need a large sample size due to the inherent variability of the cognitive assessment made by the Morris water maze or other behavioral experiments on mice. In brief, if a common specific target or pathway is found, drugs aimed at this common target can be developed, which will be of great practical value for the early diagnosis and prevention of dementia and presbycusis.
AUTHOR CONTRIBUTIONS
MX and YShu: substantial contributions to the conception or design of the work. YShe, BY, PC, QW, and CF: drafting the work or revising it critically for important intellectual content agreement to be accountable for all aspects of the work in ensuring that questions related to the accuracy or integrity of any part of the work are appropriately investigated and resolved. MX: final approval of the version to be published. | 9,853.2 | 2018-06-08T00:00:00.000 | [
"Biology"
] |
Evaluation of Montenegrin Seafarers’ Awareness of Cyber Security
e-mail<EMAIL_ADDRESS>Topics on maritime cyber security have undoubtedly been attracting great public attention in recent days. The reasons are rapidly evolving computing technologies and digitalization in maritime sector. A successful cyber-attack may have catastrophic consequences and a harmful impact on people, properties or marine environment. In addition to numerous factors that pave the way for a successful cyber-attack on ships, human errors are also in the limelight as they are notorious sources of cyber-attacks today. In this research paper, the authors examine Montenegrin seafarers’ level of familiarisation with current cybersecurity risks by conducting a structured survey questionnaire. After thoroughly analysing the collected answers, the authors realise that the respondents have an insufficient level of cybersecurity knowledge and awareness. Lastly, using the quantitative risk assessment method, the authors propose the best practices for maritime cyber security in the form of implementation of mandatory training course. Evaluation of Montenegrin Seafarers’ Awareness of Cyber Security
INTRODUCTION
A successful cyber-attack may be an important issue from the safety, environmental, and commercial standpoints. Cyber security at sea is largely related to critical infrastructures and, therefore, there is an urgent need to do the revaluation of the current awareness and preparedness of crews to adequately respond to maritime cyber risks.
"Maritime cyber risk refers to a measure of the extent to which a technology asset is threatened by a potential circumstance or event, which may result in shipping-related operational, safety, or security failures as a consequence of information or systems being corrupted, lost or compromised" (International Maritime Organization, 2017a). As a matter of fact, modern vessels rely heavily on remote monitoring and automation that can provide porous holes to hackers and cybercriminals, resulting in a compromise of vessel's key components such as ECDIS, VDR, RADAR/ARPA, GNSS, ballast/cargo/engine control systems, which are operated and controlled by the crew. Skills of a crew define how efficiently the systems will work (Yousefi and Seyedjavadin, 2012).
In order to mitigate cyber-security risks and reduce the level of their human dependency, several leading maritime organizations such as e.g. IMO, BIMCO, International Chamber of Shipping developed a set of guidelines. Their purpose is to assist shipowners and vessel operators in reducing the chance of a successful cyber incident, and to recover from it.
Various internationally required training courses, such as Security Awareness Training for all Seafarers or the Marine Environmental Awareness, have already been established. IMO "encourages Administrations to ensure that cyber risks are appropriately addressed in safety management systems no later than the first annual verification of the company's Document of Compliance after 1 January 2021" (International Maritime Organization, 2017b: 1) which is a great step forward towards achieving global shipping goals.
This research paper sheds a light on why the cyber risks at sea are still not adequately treated from the seafarer-education point of view, even after some very significant events such as hacking of Maersk's assets.
In fact, the hacking of Maersk occurred back in June 2017. At that time, due to successful NotPetya malware attack, the giant company lost between USD 250-300 million, and was forced to reinstall more than 4,000 servers and 45,000 PCs (A.P. Moller -Maersk, 2017;Cimpanu, 2018). Up to date, this has been the most serious attack of its kind, once again confirming that shipping companies are not prepared to respond to cyber risks adequately.
The largest number of safety incidents at sea occur due to human error (Yousefi and Seyedjavadin, 2012). It is no different in cyber security either. Overall situational awareness of the navigator while performing his duties on the navigation bridge consists of spatial, task, and system awareness, including cyber security awareness as well (Hareide et al., 2018). Research (Svilicic, Rudan, et al., 2019) explores cyber-security threatening Integrated Navigational System, stating that cyber-security awareness of crew is satisfying. However, another study (Svilicic, Kamahara, et al., 2019: 10) states that "crew is not familiar with cybersecurity policies, procedures and agreements, and practice insufficient cyber hygiene". More articles assess the IT infrastructure and onboard policies related to cyber-security protection, but only few of them are aimed at defining seafarers' level of awareness and knowledge of cyber threats (Bolat, Yüksel and Yüksel, 2016).
Is the awareness of Montenegrin seafarers of cyber-security high enough to make them a reliable part of the defensive shield to prevent malicious attacks on board vessel? To that end, this research paper carries out the analysis of seafarers' awareness and their knowledge of basic cyber-security aspects, and further weights the findings on a risk scale.
This paper is organised as follows: Section 2 provides an overview of common cyber-security threats at sea and users' best practices. Section 3 deals with current education process of seafarers in Montenegro. Section 4 explains the method of obtaining survey responses, which are assessed in Section 5. Section 6 elaborates the problem solutions. The findings are discussed in Section 7.
COMMON CYBER-SECURITY THREATS TO SHIPS AND USER BEST PRACTICES
There is a difference between general maritime security and maritime cyber security. While the topic of the former has been widely explored since the implementation of ISPS Code in 2003, the latter requires further attention.
Various studies have been done to clarify and explore cybersecurity risks and threats on vessels. The most important ones are: Witherby Publishing Group, BIMCO, and the International Chamber of Shipping (ICS), 2019. a.
Malware -a malicious piece of code that is utilised by cyber pests to carry out a cyber-attack. The example of malware incorporates viruses, worms, Trojan horses, ransomware, spyware, bots, etc. The malware can steal, delete, encrypt or damage sensitive data without knowledge of the victim. "Malware often infects ship's computers through the crew's use of memory sticks". (Riviera, 2020) Social engineering -technique that manipulates human psychology to get sensitive data. The victim makes mistakes that lead to data breaches. According to the Korean Register of Shipping, "social engineering means to secure access rights to systems, data, and buildings by exploiting human psychology instead of a technical hacking technique to steal into the system". There are different types of social engineering such as: a. Phishing -combines social engineering and technical methods to trick victims into divulging sensitive information such as identity and financial-related data or anything else that attackers perceive to have value (Furnell, Millet and Papadaki, 2019). Successful phishing attack can create extreme harm, e.g. in case of steeling sensitive information about the ship or itinerary details; b. Spear phishing -yet another form of phishing. Clicking on the link may cause installation of malicious software, trackers, loss of credentials, personal data or valuable shipping details. Spear phishing is sophisticated and difficult to detect; c. The so-called e-mail spoofing, still a surprisingly easy technique used for distribution of forged electronic documents that attempt to mislead the recipient about the origin of the message (Hu, Peng and Wang, 2018).
Following of e-mail instructions or requests may lead to the loss of sensitive information, e.g. ship's schedule, data on nationality of the crew, etc.; c.
Distributed denial of service (DDoS) attack is a kind of cooperative attack model where attackers use many machines to simultaneously launch DoS attacks causing the target's resources or network band-width to become exhausted or to collapse (Li et al., 2018). On board ship, it can lead to failure of navigational, engineering, and other system.
By conducting a literature research, two main types of best practices for reducing cyber-threat at sea are identified. The first is related to the network arrangement and implementation of various software and hardware solutions, while the second is focused on asset management and user best practices. For the purpose of this research, the authors have identified widely accepted cyber-security best practices whose level of success depends on user behaviour: a.
Use a strong password -Using a strong password can create the main barrier against cybercriminals. A weak password can be guessed within hours. Hackers compromise seafarers' passwords using various techniques such as a Brute-force attack, dictionary attack, and phishing attack. This type of attack can have devastating consequences; b.
Stay vigilant against phishing emails -The seafarer should avoid clicking on any attachment or link from suspicious emails, especially when working on a ship's system or a network; c.
Avoid using removable media -Removable media such as flash drive or smartphone memory card are vulnerable devices and can pose a serious challenge to ship's systems or/and network. Therefore, a seafarer must avoid using flash drives. They should save essential ship-related documents into the cloud drive or a soft copy into a secure personal computer or a laptop; d.
Stay vigilant against SMS attacks -Seafarers often prefer using SIM cards that offer cheap rates and data plans. Today's hackers better understand human psychology and know how to manipulate it. To this end, they send a phishing SMS with a link that involves the cheapest offers on calling and data plans. As soon as the seafarer opens the link, malware is installed on his/ her phone. To avoid this nightmare, the seafarer must disregard such SMS or avoid opening unknown links inside it; e.
Avoid using free Wi-Fi -Free offers and gifts often grab everyone's attention, but they can prove detrimental to Seafarer's digital property. Threat actors often cleverly provide free Wi-Fi at ports or its suburbs. The seafarer must not access a free public Wi-Fi hotspot and must avoid putting sensitive credentials; f.
Patching -All the ship's systems should be regularly patched and updated. A patch can fix a security vulnerability and bugs in the software application as well as improve its performance. For example, if an ECDIS is not stable upon installation of new charts, a new patch can resolve the issue.
CURRENT EDUCATION PROCESS RELATED TO CYBER SECURITY AT SEA
How does the current educational process in Montenegro look like and is it good enough to suit the needs of today's market?
The education of seafarers in Montenegro is organised in two levels. The first is secondary education, which lasts for 4 years. Upon completion of Maritime High School, a person can choose between two paths -joining a vessel and starting a professional career or enrolling one of the accredited Maritime Faculties in order to get a higher education degree. Enrolment to a university study programme is allowed to anyone who has completed secondary education, even if it is not through Maritime High School. Upon completion of 3-year studies, students get the Bachelor's degree and are allowed to start their seafarer career.
In 2010, Maritime High School in Kotor started carrying out re-qualification courses for all those who had previously obtained a non-maritime high school diploma. Their purpose is to offer an alternative to people who are not interested in higher education and at the same time have no will to study the complete maritime-high-school programme for another 4 years. The re-qualification course plan has been done in accordance Exploring the current curricula of the above-mentioned educational institutions in Montenegro, one finds that each of them is exploring specific IT fields, most notably the application of software programmes from MS Office -Word and Excel. The use of the Internet and modern maritime technologies such as Radio Frequency Identification (RFID) is addressed to a lesser extent. The Faculty of Maritime Studies of Kotor goes a step further by providing students with an education in the basics of computer networks and network protocols. A deeper study of computer networks as well as their protection has not been addressed so far.
The following sections present a survey done among active seafarers in order to scale their level of awareness of cyber security.
SURVEY METHOD
Taking into consideration common cyber threats and best practices presented in Section 2 of this paper, the authors created a structured survey questionnaire. Its purpose is to find out the level of Montenegrin seafarers' awareness and their potential ability to adequately respond to cyber threats.
The total number of active seafarers licensed in Montenegro is 3,000 (official data not published). The research population consists of 429 participants sailing in the rank of deck/engine officer or Master on ocean-going vessels operated by various worldwide reputable companies, including Mediterranean Shipping Company -MSC, Mitsui Ocean Line -MOL, Eastern Mediterranean Maritime, Dabinović, Reederei Nord, Crnogorska Plovidba, Bernhard Schulte Ship Management, CMA CGM, Subsea 7, and others. Even though all the previously mentioned companies employ multinational crews, the conducted survey was limited to seafarers of Montenegrin nationality whose names are undisclosed due to privacy.
The survey questionnaire consists of a total of 18 questions which are presented in Table 3. They are structured in a comprehensive way to enable quantitative research as plausible and affordable method for gathering information from seafarers. The respondents were asked to choose only one answer for each question.
RISK ASSESSMENT AND SURVEY RESULTS
To carry out risk assessment, it is necessary to define the key terms at the very beginning: risk, hazard, harm (impact), likelihood, severity, and risk assessment.
There are several definitions of risk. A commonly-used glossary (Committee on Foundations of Risk Analysis, 2015) offers 7 definitions of risk while ISO (ISO, 2009) shortly defines it as the "effect of uncertainty on objectives". "Information security risk comprises the impacts on an organization and its stakeholders that could occur due to the threats and vulnerabilities associated with the operation and use of information systems and the environments in which those systems operate" (Gantz and Philpott, 2013).
"A hazard is a source of potential injury, harm or damage. It may come from many sources, e.g. situations, the environment or a human element. " (Maritime and Coastguard Agency, 2019, p. 37) Harm or impact can be defined as the degree of damage or harm caused to the organisation or an asset.
The likelihood of occurrence is the probability that a cybercriminal will initiate a threat or the probability that a threat could successfully exploit the given vulnerability (ISO, 2009).
Both likelihood and impact can be viewed in either objective or subjective terms. In an objective expression, likelihood and impact could be expressed in terms of numerical values. On the other hand, subjectively, both elements are termed qualitatively or utilizing a range of descriptions on a scale. with IMO model courses 7.03 -Officer in Charge of Navigational Watch and 7.04 -Officer in Charge of Engineering Watch, and are popular among elder population.
To better demonstrate the official educational path of seafarers in Montenegro, the following scheme has been created ( Figure 1). Of course, seafarer education process does not end on finishing Maritime Faculty studies, Maritime High School or a requalification course, or on joining a vessel for the first time. Table 2.
Numerical values of answer choices.
Severity is the amount of damage that a hazard could create. For example, the severity of harm can be slight, moderate or extreme.
Risk assessment is a systematic process of determining the number of hazards or threats that could occur in a given amount of time to your computer systems and networks (Prowse, 2017). "The purpose of risk assessment is primarily to support decisionmaking, including decisions on risk-reducing measures in the context of a structured, systematic and documented process" (Vinnem and Røed, 2020, p. 78). There are two types of risk assessment -i.e. "Quantitative Risk Assessment" and "Qualitative Risk Assessment. " Quantitative risk assessment is a systematic risk-analysis technique used to quantify the risks associated with the IT infrastructure of an organization. It helps in understanding the exposure to risk of the IT environment, employees (or seafarers), corporate assets and its reputation. As said before, this technique involves numerical values. Though Quantitative Risk Assessment is easier, cheaper, and quicker, it cannot give a total asset value for a potential monetary loss. For instance, using this approach we can assign the ranges from 1 to 50 or 1 to 100. If the number is high, the likelihood of occurrence is high. For example, the computer having no firewall or antivirus programme has a high probability of risk.
"Risk analysis methods that use intensive quantitative measures are not suitable for today's information security risk analysis" (Karabacak and Sogukpinar, 2005, p. 148). However, to measure cyber security awareness of Montenegrin seafarers, the authors implemented ISRAM (Karabacak and Sogukpinar, 2005) quantitative risk assessment method as the second most useful in comparison with SANS, OA, Mehari, COBRA and FAIR (Svensson, 2017).
The risk model of ISRAM is based on the following formula: Risk: a single numeric value for representing the risk.
Note: All the survey participants answered all the questions from the questionnaire. Therefore, m is equal to n.
On completion of the questionnaire, but before conducting the survey, the authors "weighted" each question to scale their importance in assessing final risk. In other words, not all questions contribute equally to the conclusion of this research. Weight scale is shown in Table 1 for both probability and consequence of cyber-attack.
Weight value
Probability of occurrence Seriousness of consequence 0 Answer has no effect on probability and/ or consequence of cyber accident.
1 Answer is slightly effective to probability and/or consequence of cyber accident.
2 Answer is considerably effective to probability and/or consequence of cyber accident.
3 Answer is highly effective to probability and/or consequence of cyber accident. 4 Answer is extremely effective to probability and/or consequence of cyber accident.
After designation of answer choices, they are converted into numerical values as shown in Table 2 in order to scale probability and/or consequence of potential cyber accident.
Further in-depth scaling of questionnaire with probability and consequence weights included is shown in Table 3. Q16 Patching, updating and maintaining of ship's navigation system (e.g. ECDIS) is always crucial. Does your company have these security controls in its cyber-risk assessment plan? P=2 ; C=3 Q17 Is it true that a cyber-incident can go unnoticed for a substantial period and does not have to involve an obvious system fault or alarming ransomware messages? P=2 ; C=3 Q18 What is the typical sign that your vessel's IT/OT infrastructure is cyber-attacked? P=2 ; C=3 System is slow or unresponsive / 0 System displays warnings and alarms to inform the user about an on-going cyber-attack / 4 I do not know / 4 The minimum and maximum probability of cyber incident can be scaled based on survey results by using the equation (2): Calculations are presented in Table 4, where possible survey values are grouped evenly and scaled to represent the probability level of risk parameter. Table 4 is the risk table constructed for the probability of cyber-security incident parameter. As per Table 4, maximum possible value for survey result is 136, while minimum value is 5. For the purpose of this research, the interval of 'very high probability' is set to 27, while for other scales it is set to 25.
Using the same principle and replacing i with j in equation (2), the authors obtained the minimum and maximum values of the survey output to measure the consequences of cyber incident ( Quantitative risk matrix used for this research is presented in Table 6. It is a modified version of the risk matrix which is frequently seen on board merchant vessels and is widely used for risk assessment of daily tasks (Maritime and Coastguard Agency, 2019). Multiplying quantitative values of probability and consequences, the final value of the risk is obtained. Once the previous steps were completed, questions were distributed to 638 people who are active seafarers. Out of that number, 429 people fully responded to the questionnaire. Due to space constraints, Table 7 represents an extract of all the collected data, with average calculated probability [T1] and consequences [T2] of risk.
Likelihood of harm
Calculated risk based on the conducted-survey questionnaire, by application of fundamental risk equation (1) is 11.18, which can be described as medium level risk.
Respondent # m [m=n]
Probability of cyber incident ∑ i w i p i , where i = 429
Proposed model course for cyber-security awareness.
SOLUTIONS
Maritime industry is being rapidly digitalised, and IT is playing a crucial role in this regard. Before knowing how to prevent cyber-attacks, it is essential to know how these attacks are detected. Typically, seafarers are unaware of the attack and remain oblivious until a real loss occurs. It is indispensable for seafarers not only to adopt and understand new technologies, but also to keep themselves abreast of threats and attacks in the face of the ship's IT infrastructure.
Based on the conducted quantitative survey and ISRAM risk assessment methodology, authors measured the risk level of cyber-security awareness of Montenegrin seafarers. Rated as a medium-level risk, it can be treated as a clear indicator of necessity of urgent actions.
The human factor is always crucial when it comes to the cyber security of a ship, and this is also an important subject of this research paper. To that end, the authors proposed a model of the training course that should be set as mandatory for all crewmembers. The model course is presented in Table 8. The proposed training course should be set mandatory for all crewmembers, and it should continue in form of refresh courses on a regular 5-year basis. Implementation into the existing IMO model course 3.27 -Security awareness training for all seafarers, is also possible.
As per Security familiarisation training is also essential before joining the ship's duties. The Shipboard familiarisation checklist should be expanded to include cyber-security related training, which should be performed by the Ship Security Officer or an equally qualified seafarer. Familiarisation process should be adequately structured to guide a newly joining seafarer how to report a security incident, how to act in IT security-related emergencies and to explain which security solution is required in the event of a cyber-security incident.
CONCLUSION
In the world of digital warfare, the global shipping community including vessels, ports, terminals and various other facilities are relying heavily on the Internet to establish connectivity. Automated equipment, GNSS, ECDIS, AIS, engine/ ballast/cargo control systems, and consignment tracking systems are just some of the items dependent on adequate cyber security.
Policies and procedures on board ships should be structured and planned, accompanied by an appropriate IT infrastructure including firewalls, anti-malwares, etc. Ship's IT infrastructure is vulnerable to cyber-attacks, and human error can play its part in this regard. Therefore, achieving the overall cyber security of the ship is out of the question without a proper and effective training of seafarers. The conducted risk assessment based on survey questionnaire implicitly shows that human resources are a hot topic in terms of cyber security on board ships.
Montenegrin seafarers are mostly novices with regard to IT and cyber security. In addition, they have not acquired any IT and cyber-security related education from shore-based institutions either. For example, neither Maritime High Schools nor Maritime Faculties in Montenegro are providing any sort of education about cyber security at sea. Thus, maritime cyber security of Montenegrin seafarers is not up to the mark and needs urgent attention.
Therefore, a holistic approach to cyber security should start with the increase of people's awareness and focusing of knowledge on the mindset with appropriate training. If their training is planned to make them aware and ready to act on any threat, there is no doubt that the overall risk will be significantly reduced.
Implementation of the authors' proposed training course would set a milestone on security at sea. The proposed model course in cyber-security awareness would help in protecting confidentiality, integrity, and accessibility of information through various measures relating to people, processes, and IT systems on board ships.
Further research should focus on developing unique teaching syllabus of cyber security that will suit the needs of both Montenegrin seafarers and their employers. | 5,496.4 | 2020-10-21T00:00:00.000 | [
"Computer Science"
] |
Fast Convergence Methods for Hyperbolic Systems of Balance Laws with Riemann Conditions
In this paper, we develop an accurate technique via the use of the Adomian decomposition method (ADM) to solve analytically a 2 × 2 systems of partial differential equation that represent balance laws of hyperbolic-elliptic type. We prove that the sequence of iteration obtained by ADM converges strongly to the exact solution by establishing a construction of fixed points. For comparison purposes, we also use the Sinc function methodology to establish a new procedure to solve numerically the same system. It is shown that approximation by Sinc function converges to the exact solution exponentially, also handles changes in type. A numerical example is presented to demonstrate the theoretical results. It is noted that the two methods show the symmetry in the approximate solution. The results obtained by both methods reveal that they are reliable and convenient for solving balance laws where the initial conditions are of the Riemann type.
Introduction
Many mathematical and practical problems in physics can be modeled by balance laws, which are systems of nonlinear partial differential equations, usually of hyperbolic type. Balance laws describe very large branches of scientific fields, specifically in fluid and gas dynamics, quantum mechanics, and astrophysics.
In this paper, we present the mathematical formulation for finding numerical solution using Sinc basis functions for 2 × 2 systems of balance laws, while for comparison purposes, we also use the Adomian series technique to find approximate analytical solutions for the same model system that which can be written as The solution U, which is a vector of two components u(x, t) that might represent the length of a wave, and another component, say v(x, t) which is the velocity of that wave. i.e., U(x, t) = [u(x, t), v(x, t)] T . On the other hand the flux function can be represented as F(U) = [ f (u, v), g(u, v)] T , and finally, the source term has the form G(U) = [h 1 (x, t), h 2 (x, t)] T . where F is a smooth function on IR 2 of U, and the derivative of F possesses strictly nonlinear or linearly degenerate eigenfunctions and eigenvalues. The system (1) is of mixed type, it is of elliptic type, for all values of u, v that satisfy ε = {(u, v) ∈ R 2 : ( f u + g v ) 2 < ( f u + g v )} and of hyperbolic type for data that lying in H = {(u, v) ∈ R 2 : ( f u + g v ) 2 > ( f u + g v )}. We consider in Equation (1) the simplest piece-wise constant initial data, i.e., for certain dispositions of the Riemann initial data For the sinc method, the basis functions on the interval (−∞, ∞) for z ∈ D E , and for the time interval (0, T) are derived from the composite translated sinc functions and, π((z−kh)/h) , z = kh 1, z = kh, k = 0, ∓1, ∓2, . . .
Let f be a function defined on IR, then for h > 0 we define the Whittaker cardinal series The properties of (6) were studied and surveyed in depth in Stenger and in Lund [15][16][17]. They are based on the infinite strip D d defined in C as To approximate by the Sinc method, we will refer to an important class of functions, see [17] called Play-Wiener class, denoted by W(π/h) which is the family of all analytic functions that are in C, and satisfy some decaying conditions. For a function f and its derivative f to be in the class W(π/h), the kth derivative of the function f can be approximated by i−n denotes the kth derivative of the cardinal function. Now we present a fundamental definition that has an important role in the development of Sinc methods on arc Γ in C. For the simply connected domain D, let φ be a map that take D onto D d with φ(a) = −∞ and φ(b) = ∞. If the inverse map of φ is denoted by ψ, define When we evaluate a function and its derivative with respect to φ, we obtain Since we will replace the first derivative by it matrix approximations, here we define the m × m matrix I (1) whose in − th entry is given by δ (1) i−n . An important class of functions denoted by L α (D) play crucial rule in approximating derivatives will be needed in next formula (see, [17]). Let F ∈ L α (D), α is a positive constant, then for h = πd (αN) we have An application to the above formula, and in order to approximate the first derivative with respect to x for the function u(x, t) in the domain (−∞, ∞) leads to the approximation where With m x = 2N + 1 and the skew-symmetric matrix I (1) m x is defined by (9). Finally, we shall indicate a general formula for approximating the integral ν a F(t)dt, ν ∈ Γ. If F ∈ W(π/h), then for all k ∈ Z, To simplify the situation, we define the new matrix I (−1) , as, for N positive integer, define the (2N + 1) × (2N + 1) matrix I (−1) by
The Sinc-Galerkin Method: Balance Laws
Among the benefits of the Sinc method, is that the approximate solution is approaching the exact solution exponentially, as well as the Sinc method, it carefully handles systems in which the state changes from elliptic to hyperbolic and vis versa, as this depends on the change of the eigen-values of the system. To find an approximate solution for the balance law appeared in (1), without loss of generality, we begin with a 1-dimensional balance law defined in the domain R × (0, T 0 ) where the initial condition is taken to be Suppose k(x) ∈ L α (D), and for the step-size h = πd αN . Integrate Equation (12) with respect to time, and collocation with respect to the variable x will produce a system of Volterra integral equations of the form where the column vector u(t) represent u(t) = (u −N x (t), . . . , u N x (t)) T , here we use the notation u i (t) = u i (x i , t), k = (k(z −N x ), . . . , k(z N x )) T and H = (H(z −N x ), . . . , H(z N x )) T , while the square matrix A = −1 m x represents the discrete Sinc approximation for the first derivative. Subsequently an application of the use of the indefinite integration formula with respect to time variable t, and since our t ∈ (0, T 0 , then the domain . . , N t , to be the basis function for the interval (0, T 0 ). Now, we define a matrix E = h t I as defined in (11), with m t = 2N t + 1. Then the solution of Equation (12) is in matrix form is given by the rectangular m x × m t matrix U = [u ij ]: we use • to represents the Hadamard multiplication. For convergence purposes for the obtained solution we state two Theorems, in which the proof can be done by just resemble the theorems in [16].
Theorem 1 ([16]).
Let the matrix U be as defined in (15). Then for the number of points in both space and time N x , N t to be bigger than some constant that depends on d, α, we have To guarantee that our approximate solution as obtained in the discrete system (15) converges uniquely to the true solution, we state the following theorem.
Theorem 2 ( [16,18]). For any positive constant R, we can find another positive constant T 0 so that whenever U 1 − U 0 < R 2 , then the solution of (15) is unique. Also, the iteration system leads to a solution that converges to the unique solution.
Next, we state and prove a fact regarding the behavior of the solution, whenever the initial condition satisfies a decay condition. Consider the initial value problem for a balance law of the form with initial condition u(x, 0) = u 0 (x), x ∈ IR. We need to assume some decaying conditions on the right hand side of Equation (16).
Definition 1.
We say a function u(x, t) decays exponentially in x and uniformly for t ∈ (0, T) if |u(x, t)| ≤ C exp(−α|x|) for some positive constants al pha, C independent of t.
, and can be written as u 0 (x) = sech (αx)g(x), where the function g(x) is analytic in the strip D d and grows no more than a polynomial in the strip D d , and suppose the source term h(x, t, u) and its derivatives with respect to x, y, u are all exponentially decaying in x and uniformly for small t. Then the solution u(x, t) of (16) is in the class L α (D).
Proof. In order to find the solution of (16), we consider the system of first order ODEs.
such that The system (17) with condition (18) implies that However, h(x, t, u) and its derivatives with respect to the variables x, and t are all exponentially decaying in x and uniformly for small t. This implies that is Lipschitz in the two variables u, x. In addition, so, by the Picard theorem, the solution u(x, t) is unique.
Treatment of Non-Zero Boundary Conditions
The obtained solution in (15) was valid if the initial condition satisfy some certain decaying conditions, i.e, the initial condition has to belong to the class L α . In Section 5, we will discuss problems which have non-zero boundary conditions with the Riemann-type initial condition, which means that the initial condition is no longer in L α . To illustrate the situation, we consider the problem in (1) with the initial condition (3), and since the Sinc functions composed with the various conformal maps, S(m, h x ) • φ, are zero at the end-points of the domain of the problem, and since the boundary conditions are non-homogeneous Dirichlet conditions, then the change of variables (similarly for v(x, t)), will convert system (1) to another one with homogeneous Dirichlet conditions, with smooth initial condition. We substitute the transformation (19) into Equation (1) that leads to a new balance law for the new unknownũ(x, t) with an initial condition that is in the class L α for |x| large. Therefore, to find the solution forũ(x, t) we can proceed as mentioned above.
The Adomian Decomposition Method (ADM)
A lot of published research papers talked about the method of Adomian, and presented it in a detailed way. Here, for the sake of non-repetition, we will suffice to mention the references that contributed to the development of the method. In addition, because of its abundance and where we cannot limit it, we mention the following [19][20][21][22]. We start by rewriting the system in (1) as follows: (20) or, we may express Equations (20) using linear differential operators as In the above, the operators L t , L x are defined to be L t = ∂ ∂t and L x = ∂ ∂x . Performing the inverse operator L −1 t = t 0 . . . dt PDE system in (21) yields where The ADM [19,20] suggests that our approximate solution for unknown functions u(x, t) and v(x, t) are given by the infinite series while the four nonlinear operators φ i (u, v), i = 1, 2, 3, 4 can be expressed as an infinite series of Adomian polynomials as φ ( u, v) = ∑ ∞ n=0 A ni , where A ni are what is called Adomian's polynomials that can be computed according to specific formula set in [20]. Here, the general formulas for Adomian's polynomials A ni , i = 1, 2, 3, 4, has the form These formulas are smoothly found using coding programs to get as many terms as we want. Plug the nonlinear system (23) into Equation (22) and we arrive at the recursive relation The construction of the solutions of u(x, t) and v(x, t) are given by Regarding the convergence of the above two series, there are previous studies that dealt with the topic in detail and in general, we mention [21].
Convergence of the ADM Approximation
Here we give a proof to show that our obtained solution using ADM converges to the exacts solution via the use of the fixed-point theory. The partial differential equation that we are encountered to solve is which can be written as where H is the Hilbert space, and Comparing terms, we end up with the algorithm u 0 = F(x, t), u n+1 = A n (u 0 , u 1 , . . . , u n ). Then the approximate solution is given by a finite sum of the series S n = ∑ n n=1 u n = u 1 + u 2 + · · · + u n . If S 0 = 0, S n+1 = N(S n + u 0 ) = ∑ n k=0 A k . If the limit exist, then S = lim n→S n is the solution of the fixed-point functional equation S = N(u 0 + S) ∈ H. Need to show that the sequence {S n } is Cauchy. The distance between two consecutive iterations S n and S n+1 is a ball S r of radius r is given by To prove that {S n } is Cauchy sequence, then for all positive integers m, n where m > n, we have Choose r ∈ (0, 1) such that, S m − S n can be made arbitrarily small provided n is large enough for all m, n ∈ IN. With this choice of r and u 0 , u 1 , . . . we see that all iterations will remain in the ball There is an integer N such that u m − u n < , ∀n > N, for any m. So the sequence is Cauchy in the Hilbert space. Hence converges to some S ∈ H., i.e., lim n→∞ S n = S, where S = ∑ ∞ n=0 u n . Solving Equation (26) is the same as solving the functional equation N(S + u 0 ) = S, by assuming that N is a continuous operator, we get When viewing numerical calculations in the last section, and for purposes of comparison we will calculate what is called the order of convergence, we will provide the following definition Definition 2. A sequence S n is said to converge to S of order p ∈ IN, if there is a positive constant α such that Let us examine the order of convergence for the Cauchy sequence S n . We write Taylor's series for N(S n + u 0 ) about (S + u 0 ) as However, it is known that N(S + u 0 ) = S and N(S n + u 0 ) = S n+1 , we may rewrite the above equation as Assume our operator N ∈ C p [a, b] with the property that N (m) (S + u 0 ) = 0, for all m = 1, 2, . . . , p − 1, while it satisfy N p (S + u 0 ) = 0, therefore we have divide by (S n − S) p and take limit as n → ∞, we obtain Using the fact that S n is Cauchy that converges to S, the above equation reduces to Therefore, the order of convergence of S n is p.
For the purposes of calculating the order of convergence, we assume that the first four iterations of the computed solution are given by x k−2 , x k−1 , x k and x k+1 so the following formula will be used to calculate the order of convergence As mentioned in [23], the first study comparing ADM and Picard's method was in 1987 by Rach [24]. The study was done by reviewing several examples. Rach showed that solving differential equations of linear type using ADM was equivalent to the classical method of successive approximations, which is known as Picard iteration. Also it was shown that the Picard method gives more accurate solutions, but it needs more time than ADM; this happens because the main advantage of Adomian's method over the Picard method is the ease of computation of successive terms.
Applications: Riemann Type
To investigate the accuracy of the ADM compared with the Sinc method, we choose examples with known solutions that allows for a more complete error analysis. So a 2 × 2 system of balance laws with Riemann type conditions is tested numerically by using the Sinc function methodology, and for comparison purposes, we also solve the same model by ADM. The example reported here is selected to show the convergence of the two schemes. Consider the system where T is some small constant. In addition, with flux function F; IR 2 → IR 2 , such that with the approximations to Riemann type condition given by The left and right boundary conditions are .
In Section 3 when we construct the approximate solution by means of the Sinc methodology, it was necessary that the initial condition belonged to the class of functions L α (D) previously mentioned in Section 3, but we know that the Riemann condition does not belong to this family, so a transformation had to be made to make the primitive condition belong to the family. Since the boundary conditions are non-homogeneous, then the transformatioñ will convert the given boundary conditions to homogeneous conditions, provided that |x| → ∞. Now, after substituting the transformation (29) into Equation (28), we get a new system with the unknowñ U(x, t) = (ũ(x, t),ṽ(x, t)) T given by where W(x, 0) = (w 1 (x, 0), w 2 (x, 0)) T . Equation (30) can be solved forŨ(x, t), and then, using the transformation in (29), for U(x, t). To proceed, since the space and the time domains are (−∞, ∞) and (0, T), respectively, choose the conformal maps φ(x) = x and Υ(t) = log( t T−t ). We solve the above system using ADM subject to the initial conditions u(x, 0), v(x, 0). Integrate the system in (28) with respect to t and using the initial conditions u(x, 0), v(x, 0), we arrive at and, v( In the above two integrals, set u(x, t) = ∑ ∞ n=0 u n (x, t), v(x, t) = ∑ ∞ n=0 v n (x, t), and for the nonlinear terms, we find Adomian polynomials as The first few terms of Adomian polynomials are given by Substitute the above infinite sum for Adomian polynomials, together with the approximate solutions into Equations (32) and (33), exactly as mentioned in Section 5. Balancing terms in Equations we obtain an approximate solution of the form Table 1 shows that the method converges for T = 2 using Sinc methodology, where the second column reports the supremum norm of the error between the exact solution and the Sinc approximate u−solution, while the third column reports the error between the exact solution and the Sinc approximate v−solution. The error in the approximate u-solution using ADM is reported in Figure 1. Figure 2 reported the v(x)-solution using ADM for different values of time t. The error in approximating u, v show that most of the error concentrates at origin, which caused by approximating Riemann type condition by a smooth function. Through our study, and upon consulting the Tables 1 and 2, we are able to find the values for the order of convergence using Equation (27), it was found that the order of convergence for the ADM is q = 1.0007, while it is approximately q = 2.843 using the Sinc method. This is in accordance with the theories formulated in Section 4, which show that the Sinc methodology is faster, better, and more comprehensive in finding an approximate solution. Because of the nature of the issues under discussion, it is noted that the two methods show the symmetry of the approximate solution. A quick look at in calculated approximate solutions for both u(x, t) and v(x, t) shows the symmetry of the solutions where the error value at the point (x, t) is the same as at (−x, t).
The main goal of this research paper is to compare the two methods together, but if the reader wants more ways to solve systems from the same context, we cite the following [22]. It was shown in [25] that traditional methods, like finite difference method is unstable when the system is of mixed-type when Riemann conditions are involved, which means that the two methods used in this research are an improvement over the previous traditional methods. To present a physical and engineering application for the purposes of explaining the importance of the topic, please review the research paper [18] where the p−system was solved. Finally, what we intend to refer to in a future research in the hope of discussing it from all sides is the importance of the Telegraph system. As an illustration, we shall consider conductors (such as telephone wires or submarine cables) in which the current I = I(x, t) may leak to ground. The resultant decrease in current is governed by I x = −GV − CV t , where V = V(x, t) is voltage, G = G(x) is the conductance to ground, and C = C(x) is the capacitance with the ground. The change in voltage is governed by V x = −RI − LI t , where R = R(x) is the resistance and L = L(x) is the inductance of the cable. | 5,159 | 2020-05-06T00:00:00.000 | [
"Mathematics"
] |
Immobilization of Glucanobacter xylinum onto natural polymers to enhance the 1 bacterial cellulose productivity 2
11 Bacterial cellulose (BC) has profound applications in different sectors of 12 biotechnology due to its unique properties preferring it about plant cellulose. 13 Although this polymer is extremely important in various applications, many problems 14 still hinder the sustainable production in terms of increasing productivity and low-cost 15 production. In order to overcome these problems, this study will focuses on the 16 continuous production of cellulose using immobilized Glucanobacter xylinum cells 17 onto Sugar Cane Bagasse (SCB) and Ca-alginate beads. Comparatively, adsorption of 18 Glucanobacter xylinum cells to the cavum of stalk cells of SCB could be efficiently 19 stable while, entrapment of cells onto Ca-alginate has drawback observed by the rapid 20 disruption and instability of the beads in the Potato Peel Waste (PPW) culture 21 medium. Our findings demonstrate that a combination between alternative low-cost 22 medium with continuous production mode by immobilization onto inexpensive 23 natural polymer can promote a sustainable bioprocess and reduction the production 24 cost. 25 Utilization of waste from the food industry as raw materials for both immobilized the bacterial cells and prepared the culture medium promotes economic advantages because they reduce environmental pollution and stimulate new research for science sustainability. The observed study was carried out to produce bacterial cellulose via immobilization onto fibrous and non fibrous bio-polymers. The foregoing results justify the applicability of SCB as carrier matrix for immobilization of BC in biosynthesis of cellulose from Potato Peel Waste hydrolysate culture medium. Reused immobilized biomass indicated sustained cellulose production even after 6 cycles. The instrumental analysis of BC produced from fibrous biopolymer showed excellent characters with high crystal structure and homogenous network as illustrated from SEM topography. These results demonstrate the feasibility of the proposed immobilization system to be used in future industrial BC production from low cost raw materials.
limitations, such as a high operating cost, rapid consumption of the substrate, low in 64 pH stability, and a low BC productivity. Several approaches have been suggested 65 to improve BC production efficiency, involving supplementation of the cultural 66 medium with some regulators such as ethanol or organic acids in order to inhibit the 67 accumulation of the basic metabolic byproduct (gluconic acid) and at the same time 68 stimulate the synthesis of substances necessary for the cell stabilization (Lu et al.,69 2016, Stepanov and Efremenko, 2018). However, these additivesare not suitable, due turn caused an inhibitory effect on BC biosynthesis (Morgan et al., 2014). Therefore, 76 to improve the BC yield it should be foster the key molecule c-di-GMP in the 77 metabolic process of the cells. As discussed by (Srivastava and Waters, 2012, 78 Stepanov and Efremenko, 2018) c-di-GMP is a major metabolic molecule called as 79 "quorum factor," since the highly cells concentrations was correlated with quorum 80 state. In this state the BC production was increased by the expression of "silent genes" 81 and the synthesis of exopolysaccharides with a simultaneous decrease in the rate of 82 active cell growth. Therefore, cells that produce BC should be stimulated to come 83 into a quorum state, which the cells become genetically programmed to their 84 increased population. The cell-immobilization system in case of BC producers could 85 allow obtaining highly concentrated populations of cells since BC synthesis would be 86 regulated by a quorum sensing phenomena as described before.Interestingly, the 87 immobilization of BC cells opens a way to improve cell stabilization and thus led to 88 increase BC productivity. Comparatively, the immobilized cells have various benefits 89 more than free cells in the production process, such as increased the cell population of the overall advantages of immobilization process the current reports concerning the 96 production of BC by immobilized-cell system are very rare. In this regard, PVA 97 cryogel was used for employing Komagataeibacterxylinum cells in an immobilized 98 system to increase the biosynthesis of BC (Stepanov and Efremenko, 2018). 99 Acetobacter xylinum ATCC 700178 cells was successfully immobilized on a plastic 100 composite support (PCS) to improve the BC production on the basis of polypropylene 101 (Cheng et al., 2009). However, the severe masstransfers restrictions, low mechanical 102 strength, non-biodegradability and highly toxicity of these synthetic polymers 103 displaying a big problem in the operational stability of the immobilized cells (Basak et 104 al., 2014, Nuanpeng et al., 2018). Therefore, we tried to finding out a renewable, easily 105 prepared, inexpensive, biodegradable, non-toxic, and available naturally carrier. studies have yet been established on the statistical optimization of BC production 125 using immobilized-cell system despite its high industrial applications. Therefore, we 126 investigated the enhancement of bacterial cellulose production by immobilized G. Egypt. Na-alginate purchase from molekula (U.K), Potato peel waste (PPW) was 136 resulting from potatoes processing, and collected from the disposal of free markets.
137
The bright PPW without disease symptoms were selected then washed thoroughly 138 with distilled water. All reagents, solvents, medium and its components used in this 139 study were of analytical grade.
151
Sugar cane bagasse (SCB) was obtained from a sugarcane juice local market 152 in Cairo, Egypt,after the skin and the outside fiber were removed; SCB was chopped 153 into small particles using a food processor. The chopped SCB was then dried, and 154 approximately 50 mL moisture was vaporized from 100 g raw SCB. Bagasse which 155 was obtained after drying was sieved to remove fine and larger particles. The pieces 156 of SCB were sieved to obtain particle sizes of 1 mm x1 mm x 1 mm, 2.5 mm x 2.5 157 mm x 2.5 mm, 5 mm x 5 mm x 5 mm, and 10 mm x 10 mm x 10mm. The crushed and scanning electron microscope. The samples for electron microscopy were prepared 195 according to the method described by (Yu et al., 2007). In all BC production 196 experiments, the reducing sugar concentration of in the PPW culture medium was measured by DNS according to the procedure reported in our previous work (Miller,198 1959). As well as, cell retention (Cr, CFU g -1 ) onto the SCB particle and alginate 199 beads were measured as the ratio of total number of CFU immobilized onto the carrier 200 to the carrier mass (g). Log CFU was determined as adapted by (Abdelraof et al.,201 2019a).
202
The immobilization efficiency (Yi, %) was calculated as follows: In order to study the reusability of the SCB immobilized cells, after every 254 batch BC production, the whole SCB contents of the flasks were collected aseptically 255 from the spent medium and washed three times with sterile bi-distilled water. Then,
256
SCB particles were used again separately for BC production with fresh PPW medium 257 under the same experimental conditions. This cycle was repeated for ten times to bacterial cells onto SCB particles and alginate beads were confirmed using scanning 309 electron microscope (SEM) compared with the non-inoculated matrices (Fig. 1). As 310 shown in (Fig. 1A, B), the bacterial cells were success attached in the alveolate of the 311 stalk cells of the SCB definitely, and high cell concentration was observed. measured up to 24 h. As shown from (Fig. 1E, F) (Table 1), the cell retention was increased as a SCB particle 378 size was increased. In fact, a larger SCB particle can be carrying more bacterial cells 379 than a smaller one, and that because it has more intact stalk cells (Basak et al., 2014). 406 The used instrumental tools are useful in characterization of produced BC 407 which included FTIR, XRD, SEM. The FTIR spectra are clearfield in Fig. (2A). The respectively.
432
The topography study is carried out on the produced bacterial celluloses 433 show significant differences in topography study in Fig. (3). The HS BC appears as 434 spongy-like this may be referring to low crystallinity as shown in Fig. (3A). The PPW 435 appears as dark spots cellulose with enhance in crystal appearance at The three dimensional (3D) response surface plots-generated by Minitab-17 software 507 is shown in Fig. (4), represents the relationships and effects of different experimental 508 variables (factors) on BC productivity. Best experimental variables levels for 509 maximizing BC production were predicted through analysis of these plots in 510 combination with numerical optimization for each variable and desirability analysis. Therefore, our previous studies attendance to the PPW medium is a successful 520 hydrolysate waste to regular production of BC without any influence with the pH 521 value and that due to its having high buffering capacity and also has a good impact on 522 the formation of biopolymer. According to the statistical bioprocess optimization, we 523 can noticed that the direct proportional between the sugar consumption and pH value.
524
From these results we can conclude that, immobilization of Glucanobacter xylinum 525 onto a low-cost abundant natural byproduct (i.e. SCB) and optimization the BC 526 production using the PPW hydrolysate medium opening an effective way to 527 sustainability of BC biosynthesis. cycle's number of repeated batch cellulose production by the SCB-immobilized cells 537 and the main fermentation kinetic parameters are summarized in Table 5. As can be 538 seen, reuse of the SCB particles could be exactly carried out for four sequential times 539 without any significant decrease in the operational efficiency of the BC yield. It was 540 observed that the BC production rate was initially affected at the 7 th cycle and that With respect to the kinetic studies, the SCB-immobilized cells exhibited slightly 554 higher cellulose productivity rate (0.043 g/L. h) through five repeated batch fermentation than the free cells (0.0401 g/L. h). After that, the cellulose productivity 556 was starting reduced and this might be due to the fact that the immobilized cells was 557 reduced in the carrier and that was clearly appeared in the substrate conversion rate 558 which decreased with 10%. To deeply understand these changes, it should be noted 559 that the utilization of sugars by the SCB-immobilized cells was not restricted with the 560 carrier system, suggested that the diffusion of the substrates was not prevented by the 561 carriers, which were highly porous and thus, facilitated the mass transfer of the suggesting that high sugar concentrations in the fermentation broth had no effect on 567 the bacterial growth. We propose, from these findings, that the regeneration and | 2,413.6 | 2020-12-01T00:00:00.000 | [
"Engineering",
"Biology"
] |
Atmospheric corrosion of Nylon 6,6 in Mauritius
Plastics are being more and more commonly used for outdoor applications in Mauritius. In this context, the atmospheric corrosion degradation behaviour of nylon 6,6 was observed in the Mauritian atmosphere, having a corrosivity category of C3, according to ISO 9223. The crack width, depth and extent of the cracks formed on the surface were investigated using micrometry and image analysis. The changes in the chemical composition of the nylon 6,6 were investigated using the Fourier Transform Infrared spectroscopy technique. It was observed that the nylon 6,6 surface showed a progressive discolouration and a sudden increase in the size of the cracks after around 2 years of exposure. The degree of crack formation in terms of crack width, the depth of degradation and the microhardness of the surface were evaluated and all showed significant degradation only after about 2 years of exposure. Degradation has penetrated the test material within 750 μm from the surface in 4.5 years of exposure. Using the FTIR technique, evidence of chain scission occurring among the polymer chains was observed.
INTRODUCTION
Any polymeric material which is mainly exposed outdoors, and hence to the weather conditions, will most likely deteriorate after a period of time. This process is also known as weathering. Weathering of polymers is a complex and unpredictable process because the degradation factors such as temperature, humidity or intensity of solar irradiation will depend primarily on the weather conditions. However, much emphasis needs to be put on the weathering resistance of polymeric materials as the demand for these materials is increasing day by day. Polymers deteriorate through: 1. Thermal degradation, which can be followed by thermo-oxidation [1]; 2. Ultraviolet degradation, such as photo-oxidation. This involves the formation of micro cracks and a yellowing effect on composites [2] or surface cracking on long fiber-reinforced thermoplastic composites samples [3]; 3. Hydrolytic degradation, which can cause subsequent discolouration of the material [4]; 4. Environmental stress cracking; 5. Any combination of the above degradation process; 6. Metal induced degradation [2]. Deterioration of synthetic polymers generally involves changes in the physical and visual appearance of the material and may appear in the following ways [5]: Surface embrittlement, crack formation, yellowing of surface, softening of material, discolouration of surface, charring, delamination and swelling of material. Some common properties, which are often investigated in numerous weathering studies, to determine the amount of degradation in synthetic polymers are: 1. Surface change-most oxidation processes, such as thermo-oxidation and photo-oxidation, take place readily on the surface. This usually results in the formation of cracks which propagates progressively across the surface. Yellowing of the material is also a common parameter used to gauge quickly the degree of degradation due to U.V exposure. A study carried out by Chevali et al. [3] showed that there was a significant increase in the yellowing effect of polypropylene (PP) with increasing exposure duration after 400 hours. 2. Change in the chemical structure-Several studies on weathered polymers suggested that photo-oxidation and thermo-oxidation generate carbonyl groups as degraded products during the propagation step. The formation of these particular chemical compounds is also proposed by RABEK [1]. Additionally, WOO et al. [2] identified the presence of carbonyl groups compounds in photo-degraded nanocomposites specimens using Fourier Transform Infrared spectroscopy (FTIR) and reported an increase in its concentration with increasing exposure time to ultraviolet radiation. 3. Change in molecular weight-Since most types of degradation mechanisms involve chain scission of the polymer backbone into smaller molecules, the overall molecular weight of the polymer will actually decrease. This change in molecular weight is sometimes used as an index to evaluate degradation. 4. Change in mechanical properties-Tensile strength is the most common mechanical property studied with weathered specimens because it provides important information on whether the weathered specimens will function as intended if a load is applied during their service lifetime. Qayyum and White [6] reported a considerable reduction in tensile strength on weathered polypropylene compounds. Additionally, the same authors observed in another study [7] the presence of brittle layers on polyvinyl chloride (PVC) and Nylon 6,6 weathered specimens which exhibited fracture and ultimate failure during tensile testing. In Mauritius synthetic polymers are becoming very popular for outdoor exposures. Large steel towers for fixing telecommunications antennas, for example, are now being embellished by the use of synthetic polymers to give them a more aesthetically pleasing look.
Deterioration of a polymeric component will not only affect the function and service life of the component but also its aesthetics, in some cases. One common example here is the outdoor plastic furniture which when exposed to sunlight for a very long time degrades in terms of appearance. In the long term, the most common way for polymers to degrade, when exposed outdoors, is through weathering when exposed to the atmosphere. However, degradation of polymers may also occur due to specific conditions such as heat caused by external sources.
Ito and Nagai [8] found that in the Japanese railway fields, polymeric components such as insulation roof sheets, hand straps, or air ducts for instance had an estimated service lifetime of 20 years but the polymeric products were replaced every 5-10 years for reliability. The main degradation factors responsible for the deterioration of these components under conditions in which railway services operated, were reported as mainly in terms of temperature, solar irradiation, presence of chemicals and mechanical vibrations.
In Mauritius, though synthetic polymers are being more commonly used for outdoor applications, no study has been performed to investigate how these materials degrade in the Mauritian atmosphere. In atmospheric corrosion tests, carbon steel samples were exposed outdoors, according to BS EN ISO 8565 [9]. In the exposure racks, nylon 6,6 supports were used to hold the metal samples. During the study, it was observed that the nylon supports had also weathered. Hence, in this study the atmospheric degradation in the nylon supports are investigated to get a better insight into their atmospheric degradation behaviour in the Mauritian atmosphere.
Corrosivity of the Mauritian atmosphere
Mauritius is a tropical island of 1865 km 2 . It lies about 800 kilometres east of Madagascar, as shown in Figure 1. Rainfall is abundant with an average of 2000 mm per year overall and relative humidity being frequently above 80%. Moreover, being a small island, Mauritius is surrounded by sea. These conditions are expected to influence the atmospheric corrosion rate of materials. To determine the corrosivity of the Mauritian atmosphere, low carbon steel was exposed at several sites in Mauritius, according to BS EN ISO 8565 [9]. The sites included a marine site (Belle Mare), a marine industrial site (Port Louis) and two rural sites (St Julien d'Hotman and Reduit). The results for the mass loss analysis performed according to BS 7545 [10] over 1 ½ years of exposure are shown in Figure 2 [11]. Based on the mass loss analysis, the corrosivity of the atmospheres were determined, according to ISO 9223 [12]. The results are as shown in Table 1. From these results and from results of other researches performed in Mauritius, it is expected that most sites in Mauritius, apart from the capital city (Port Louis), would fall in the corrosivity category C 3 [11]. Hence, in this study, the nylon from corrosivity category C3 sites was chosen for further experimentation and analysis.
MATERIALS AND METHODS
The nylon 6,6 test specimens used in this study were initially used to support metal plates in the exposure racks for outdoor exposure of low carbon steel samples as shown in Figure 3. They are of 32 mm diameter and 20 mm in depth. A total of 56 samples were randomly collected from the exposure racks after 1,1.5, 2, 2.5, 3, 3.5, 4, 4.5 years of exposure. All the specimens removed from category C3 environments, according to ISO 9223 [12], were considered. It should be noted that, Mauritius is a small tropical island of 1865 km 2 and the atmospheric corrosivity over the whole island does not vary much and falls in the category C3, apart from the capital city which falls in category C4 [11].
The main aim of this study is to investigate how atmospheric corrosion degradation on nylon 6,6 progresses in the island. This was performed by: 1. Examining the surface of the components to find width, depth and extent of the cracks formed. This was performed through micrometry, using the Nikon ME600 metallurgical microscope, and image analysis, using the ImageJ v1.46r software.
2. Determining the changes in the chemical composition of the nylon 6.6 with time using the FTIR technique.
3. Performing microhardness tests on the samples' surface. The crack widths on the surface of each of the test specimens selected were measured using a filar micrometer arrangement connected to the eyepiece of the Nikon Microscope ME600. The average of the 5 largest crack width observed on each nylon sample was determined.
Thickness of degraded layer
The thickness of the degraded layer across each weathered specimen was measured in order to investigate how deep the deterioration has penetrated the material. Each specimen was first sectioned. The crosssectioned surface of each test specimen was then placed under the optical microscope and the average maximum depth of the degraded layer observed was measured through micrometry.
Vickers Microhardness
Simultaneously, the variation of the hardness of the samples along their length, perpendicular to the weathered surface was measured. This was performed by cutting the samples along their diameter. The Vickers indenter was used for microhardness testing. Five measurements were taken along the depth of the cut sample at 0.5 mm distance interval, according to ATSM E384-11e1 [13]. A load of 50 gf was used because higher loads made larger indentation which was out of bounds of the field view for measurement.
Also, the variation of the hardness of the nylon 6,6 samples 0.5 mm from the degraded surface was determined for various time of exposure. This was performed to find out whether the surface hardness varies with time of exposure.
Use of FTIR
Degradation of a polymer, usually caused by thermo-oxidation and photo-oxidation during weathering, involves a chain reaction involving free radicals, which is followed by the formation of new chemical compounds as products. These degradation products commonly appear as carbonyl groups compounds. By identifying and quantifying the amount of these specific chemical groups in weathered Nylon specimens, it is possible to evaluate the degree of degradation based on this change in chemical structure. The FTIR was used in this study to investigate the change in the chemical structure of the Nylon test specimens due to weathering.
A thin surface layer was cut from each sample for this purpose. 0.0145 g of the cut sample was then dissolved in 2,2,2-Trifluoroethanol (TFE). This solution was then used in the FTIR tests.
Visual examination of degraded surfaces
An initial visual examination of the degraded surfaces at both macroscopic and microscopic level revealed a progressive surface change mainly in terms of colour and formation of cracks with increasing time of exposure.
Discolouration due to weathering
Discolouration of polymer surface was observed with increase in time of exposure. This is usually caused by photo-degradation of the surface upon UV exposure. Figure 4 shows typical progressive surface colour change and degradation observed on the weathered samples' surfaces. Yellowing or discolouration is a common phenomenon observed during weathering of polymers [3].
Formation of cracks due to weathering
Micro cracks of mean width of 5 m and length 30 m, started to appear on specimens exposed for 1.5 years. This can be observed in Figure 5, in which a typical surface of a nylon sample is shown under a magnification of 500X, using a digital microscope, showing cracks of length 34.5 m and widths 4.93 m and 5.3 m. Larger cracks on the nylon surface start to appear after 2 years of exposure, as shown in Figure 6, and the degree of crack formation suddenly increase after 2.5 years of exposure. The formation of cracks became more evident as the exposure time increased to 4.5 years. After 2 years of exposure, the average maximum crack width was found to be equal to 25 µm. An average maximum crack width of 85 µm was recorded after 4.5 years of exposure.
One possible reason for this sudden formation and increase in size of the cracks was the fact that the surfaces were continuously exposed to solar irradiation, which initiated photo-oxidation. The propagation step involved in the mechanism of photo-oxidation contributes to chain scission of the polymer chains on the surface, thus causing surface deterioration [1]. During the time that the surfaces were exposed to UV radiation, the propagation step took place continuously in the presence of oxygen at the surface. The fact that no apparent surface change was observed during the first 2 years of exposure was because generation of free radicals possibly occurred during this period (initiation step), which did not involve chain and scissions.
Average maximum crack width
The crack width of the samples was measured. The 5 maximum values were considered for analysis. Their average was determined for the various time of exposure of the samples, as shown in Figure 7. Cracks were not observed up to 1 year of exposure. As from 1.5 years of exposure, the crack width was found to increase linearly with time. The trendline, as shown in Figure 7, was found to have a coefficient of determination of 0.97. This shows that the results correlate well with the trendline.
Years of exposure under an optical microscope
The average values of the maximum depth of cracks observed on the cross-section of the samples specimens are shown in Figure 8. From Figure 9, it can be observed that the depth of cracks formed on the surface were not deep for the first 2.5 years of exposure. After this period of exposure, there was a sudden increase in the depth of cracks formed reaching around 750 µm after 4.5 years of exposure.
Changes in hardness across the specimens
No change in the hardness of nylon was observed along the cross section of the samples, for the first 2.5 years of exposure. A typical graph for one year of exposure is shown in Figure 10. After 2.5 years of exposure, the hardness of the samples near the surface (1 mm from the latter) was found to be lower than that further away from the surface. This is shown in Figure 11. The changes in hardness of the samples are near to the surface, most probably because of the degradation occurs only at the surface. Since the changes in the hardness of the samples are observed on or near the surface, the degradation of nylon outdoors can be considered to affect only the surface. This also confirms the results of the tests performed from the visual analysis. In order to evaluate how the hardness has changed over time on the nylon 6,6 surface which is the most affected region, the hardness of the specimens were measured at a depth of 0.5 mm from the surface. Figure 12 shows the results obtained for the variation of the hardness with increase in time of exposure at 0.5 mm below the degraded surface. From the graph, it can be observed that there is a gradual decrease in the hardness of nylon with increase in time of exposure. The decrease started to become more significant after 2 years of exposure. One possible explanation for this decrease was the fact that the increase in size of the cracks, as discussed previously, caused the formation of cavities within the material. The presence of cracks therefore caused the material to be less resistant to deformation and hence to a decrease in hardness.
Changes in chemical structure due to weathering
FTIR analysis was performed to monitor the change in chemical structure due to degradation. There is a possibility that photo-oxidation generates degradation products belonging to the carbonyl group. Detecting the presence of this type of chemical group may demonstrate that UV exposure and heat were the main cause of degradation. Figure 13 shows the typical absorbance FTIR spectra for samples exposed for 1 year and 4.5 years.
No evidence of the presence of the carbonyl group was found in the FTIR spectra. Carbonyl group compounds usually absorb infra-red radiation between 1680 -1750 cm -1 and only the peak due to C=O in the amide bond was observed. However, George and Browne [14] made similar observations when evaluating deterioration of Nylon 6,6 parachutes.
Nevertheless, from Figure 13, it is clear that there is a decrease in the area under the peaks of the samples exposed for 1 year and 4.5 years. The area under the peaks corresponds to the relative concentration of chemical groups and it is evident that the chemical groups found in the sample exposed over 4.5 years have decreased considerably. This decrease is generally caused by the chain scission of the polymer chains upon degradation which is common during photo-oxidation and thermo-oxidation.
Weathering of nylon compared to mild steel
In the C3 environment in Mauritius, mild steel degrades linearly with time of exposure for around 28 days [11]. For the medium and long term, it was found that the corrosion trend follows the power law.
For nylon 6,6, this is not the case. Weathering is not clearly observed in the first two years of exposure. However, after this period nylon degrades at a fast rate with linear increase in surface crack formation.
The degradation of nylon 6,6 in the Mauritian atmosphere can, thus, be represented by Figure 14. This pattern consists of an initial period of time during which no degradation is observed until propagation begins and accelerate the rate of degradation.
CONCLUSIONS
Hence, in this study, it has been observed that: 1. Most of the physical and mechanical properties investigated changed in a similar pattern with increasing exposure time. This pattern consists of an initial period of time during which no degradation is observed until propagation begins and accelerate the rate of degradation. 2. Visual examination of the surface of the nylon samples showed a progressive discolouration and a sudden increase in the size of the cracks only after 2 years of exposure. The propagation of free radicals attacking the polymers was discussed as a possible explanation for the sudden change observed. 3. The degree of crack formation in terms of crack width and crack area ratio, the depth of degradation and the microhardness of the surface were evaluated and all showed significant degradation after around 2 years of exposure. The degradation occurs on the material's surface. 4. Degradation has penetrated the test material up to 0.75 mm from the surface after 4.5 years of exposure. This is significant and should be taken into consideration when designing for the outdoor environment. 5. No carbonyl group compounds were found in the weathered samples using FTIR but evidence of chain scission occurring among the polymer chains was found after investigation of the IR spectrum. Hence, changes in nylon are not only of physical nature but also of a chemical nature. Nylon 6,6 may be a good alternative to metals in atmospheric exposures. It will resist atmospheric degradation for around two years in a C3 category environment. However, after 2 years, it will degrade gradually and in the long term, it may not necessarily be a reliable material for use as alternative to metals. | 4,501 | 2016-12-01T00:00:00.000 | [
"Materials Science"
] |
Phenol-Rich Feijoa sellowiana (Pineapple Guava) Extracts Protect Human Red Blood Cells from Mercury-Induced Cellular Toxicity
Plant polyphenols, with broadly known antioxidant properties, represent very effective agents against environmental oxidative stressors, including mercury. This heavy metal irreversibly binds thiol groups, sequestering endogenous antioxidants, such as glutathione. Increased incidence of food-derived mercury is cause for concern, given the many severe downstream effects, ranging from kidney to cardiovascular diseases. Therefore, the possible beneficial properties of Feijoa sellowiana against mercury toxicity were tested using intact human red blood cells (RBC) incubated in the presence of HgCl2. Here, we show that phenol-rich (10–200 µg/mL) extracts from the Feijoa sellowiana fruit potently protect against mercury-induced toxicity and oxidative stress. Peel and pulp extracts are both able to counteract the oxidative stress and thiol decrease induced in RBC by mercury treatment. Nonetheless, the peel extract had a greater protective effect compared to the pulp, although to a different extent for the different markers analyzed, which is at least partially due to the greater proportion and diversity of polyphenols in the peel. Furthermore, Fejioa sellowiana extracts also prevent mercury-induced morphological changes, which are known to enhance the pro-coagulant activity of these cells. These novel findings provide biochemical bases for the pharmacological use of Fejioa sellowiana-based functional foods in preventing and combating mercury-related illnesses.
Introduction
Feijoa sellowiana (Feijoa), commonly known as pineapple guava, is an evergreen shrub in the Mytraceae family that is native to South America. It is commonly cultivated in tropical and subtropical countries, such as Brazil, Uruguay, Paraguay, and Argentina, but its cultivation has been extended to other countries, including Italy. Feijoa fruit, an intensely fragrant dark green oval berry, is commonly eaten fresh or as a variety of commercially available processed foods, such as jam, ice cream, and yoghurt [1,2].
The advances in chemical composition and biological activities of different botanical parts of Feijoa have recently been summarized in a mini-review by Fan Zhu [3]. The fruit is the most utilized botanical part of Feijoa, and its nutritional value is generally defined by the presence of dietary fiber, essential amino acids, potassium, and vitamins, including vitamin C. Recent studies have added to Feijoa nutritional properties, including high folic acid content, and particularly high iodine (3 mg/100 g
Preparation of Red Blood Cells and Treatment with HgCl 2
Whole blood was obtained with informed consent from healthy volunteers at Campania University "Luigi Vanvitelli" in Naples, Italy. It was deprived of leucocytes and platelets by filtration in a nylon net and washed twice with isotonic saline solution (0.9% NaCl); the resulting intact RBC were resuspended in buffer A (5 mM Tris-HCl pH 7.4, 0.9% NaCl, 1 mM MgCl 2 , and 2.8 mM glucose) to obtain a 10% hematocrit, as previously described [41]. Intact RBC were incubated at 37 • C with 40 µM HgCl 2 for 4 h or 24 h. For experiments with Feijoa pulp and peel extracts, stock solutions, prepared in DMSO as above described, were diluted in buffer A to a final DMSO concentration of about 0.02%, in order to avoid DMSO cytotoxicity. As a control, the effect of this volume of DMSO on RBC was evaluated and found to be negligible (data not shown). RBC from each donor were used for a single assay in triplicate. Each experiment was repeated on RBC obtained from three different donors.
Hemolysis Assay
RBC hemolysis extent was determined spectrophotometrically, according to Tagliafierro et al. [41]. After simultaneous treatment with HgCl 2 and Feijoa extracts for 24 h, the reaction mixture was centrifuged at 1100× g for 5 min, and the released hemoglobin (Hb) in the supernatant was evaluated by measuring the absorption at 540 nm (A). As a positive control, packed RBC were used hemolyzed with ice-cold distilled water at 40:1 v/v, and by measuring the A540 of the supernatant obtained centrifuging the suspension at 1500× g for 10 min (B). The percentage of hemolysis was calculated as the ratio of the readings (A/B) × 100%.
Determination of Reactive Oxygen Species
ROS generation was determined using the dichlofluorescein (DCF) assay, according to Tagliafierro et al. [41]. Using this method, 250 µL of intact RBC (hematocrit 10%) were incubated with the non-polar, non-fluorescent 2 ,7 -dichlorodihydrofluorescin diacetate (DCFH-DA) at a final concentration of 10 µM for 15 min at 37 • C. After centrifuging at room temperature at 1200× g for 5 min, the supernatant was removed, and the hematocrit was re-adjusted to 10% with buffer A. RBC were then treated concurrently with HgCl 2 and Feijoa extracts in the dark for 4 h. After the incubation, 20 µL of RBC were diluted in 2 mL of water, and the fluorescence intensity of the oxidized derivative DCF was recorded (λ exc 502; λ em 520). The results were expressed as fluorescence intensity/mg of Hb.
Quantification of Intracellular Glutathione
Intracellular GSH content was determined spectrophotometrically by reacting with DTNB reagent, according to Van den Berg et al. [47]. After co-incubation with HgCl 2 and Feijoa extracts for 4 h, the samples (0.25 mL) were centrifuged, and the cells were lysed by the addition of 0.6 mL of ice-cold water. Proteins were precipitated with 0.6 mL ice-cold metaphosphoric acid solution (1.67 g metaphosphoric acid, 0.2 g EDTA, and 30 g NaCl in 100 mL of water). After incubation at 4 • C for 5 min, the protein precipitate was removed by centrifugation at 18,000× g for 10 min, and 0.45 mL of the supernatant was mixed with an equal volume of 0.3 M Na 2 HPO 4 . Then, 100 µL of DTNB solution (20 mg DTNB plus 1% of sodium citrate in 100 mL of water) was then added to the sample, and after a 10 min incubation at room temperature, the absorbance of the sample was read against the blank at 412 nm.
Estimation of Free Sulfhydryl Groups in Isolated Red Blood Cell Membranes
Free sulfhydryl groups in membrane proteins (2.5 µg/µL for each sample estimated by Bradford assay) were assayed, according to the method of Ellman [48]. To do this, 650 µL of HgCl 2 -treated RBC were washed three times in 40 volumes of 5 mM sodium phosphate buffer pH 8.8, centrifuged at 10,000× g for 20 min at 4 • C, then washed several times with the same buffer, for the complete removal of Hb. Then, 1 mL of 0.1 M Tris-HCl pH 7.5 was added to 50 µL membrane protein. The colorimetric reaction was started by adding 50 µL of 10 mM DTNB in methanol. After 15 min of incubation at room temperature, the absorbance was read against the blank at 412 nm. Blanks were run for each sample in which DTNB was not added to methanol.
Morphological Analysis of Red Blood Cells
To investigate the possible protective effect of Feijoa extracts on alterations in the Hg-induced erythrocytes' shape, we treated the cells both with 40 µM HgCl 2 as well as 20 or 80 µg/mL of Feijoa peel and pulp extracts for 4h at 37 • C. After incubation, erythrocytes were washed twice with phosphate-buffered saline pH 7.4 (PBS), and counted in a Burker chamber. The confocal laser scanning microscope analyses were performed according to Nguyen [49], with few modifications. In brief, the cells were then fixed with 2% formaldehyde for 1 h at 4 • C, then washed several times and incubated with anti-human anti glycophorin A FITC antibody for 30 min at 4 • C in the dark. Afterwards, the samples were placed on glass slides and air-dried for 1 h. The slides were dipped quickly, and gently washed stepwise with ethanol from 50% to 75%, 90%, and then 100% for dehydration. Finally, cells were fixed in 2% formaldehyde and washed three times with PBS. For confocal laser scanning microscope imaging, several randomly selected frames from each sample were captured for morphological observation and statistical strength. Excitation and emission filters were set at 488 nm and 550-600 nm, respectively.
Statistical Analysis
Data were expressed as mean ± standard error of the mean (SEM). The significance of differences was determined by one-way ANOVA followed by a post hoc Tukey's multiple comparisons test. GraphPad Prism 5 was utilized for statistical analysis.
Feijoa Peel and Pulp Extracts Protect Against Hg-Induced Hemolysis.
The rescue of Hg-induced hemolysis by Feijoa extracts was assayed separately for peel and pulp, and the results are shown in Figures 1 and 2, respectively. It can be seen that 24 h treatment of RBC with 40 µM HgCl 2 resulted in approximately 13-17% hemolysis, compared to 1-2% in negative controls, as expected based on our previous work [41]. Feijoa peel extract potently reduced Hg-induced hemolysis compared to that of the pulp, with a significant 3% drop in hemolysis at 10 µg/mL, and a steady reduction of about 1% with each doubling of peel extract ( Figure 1). Feijoa pulp treatment reduced cellular lysis in similar proportions, but the required protective extract concentration to do so was almost eight-fold greater than that of the peel extract ( Figure 2). No cytotoxic effect was found by either Feijoa extracts up to the maximum concentration utilized in this study (data not shown). concentration to do so was almost eight-fold greater than that of the peel extract ( Figure 2). No cytotoxic effect was found by either Feijoa extracts up to the maximum concentration utilized in this study (data not shown). Statistical significance was calculated with one-way ANOVA followed by Tukey's test. ** (p < 0.05) indicates a significant difference from cells lacking HgCl 2 treatment. # (p < 0.05) and ## (p < 0.01) indicate significant differences from cells lacking Feijoa extract treatment.
Feijoa Peel and Pulp Extracts Reduce Reactive Oxygen Species Production in Red Blood Cells
The fluorescence probe DCF assay elucidated the protective role of Feijoa extracts against oxidative stress in RBC, as reported in Figures 3 and 4. ROS production increased nearly two-fold in Hg-treated RBC compared to the negative control. In contrast, co-incubation with 10, 20, 40, or 80 µg/mL of both peel and pulp acetonic extracts incrementally reduced ROS production in RBC. Similar to hemolysis, Feijoa peel extract prevented ROS production more potently than the pulp, and remarkably reduced fluorescence levels by approximately 50% at 10 µg/mL, to near control values at . Statistical significance was calculated by one-way ANOVA followed by Tukey's test. ** (p < 0.05) indicates a significant difference from cells lacking HgCl 2 treatment. # (p < 0.05) and ## (p < 0.01) indicate significant differences from cells lacking Feijoa extract treatment. concentration to do so was almost eight-fold greater than that of the peel extract ( Figure 2). No cytotoxic effect was found by either Feijoa extracts up to the maximum concentration utilized in this study (data not shown). Statistical significance was calculated with one-way ANOVA followed by Tukey's test. ** (p < 0.05) indicates a significant difference from cells lacking HgCl 2 treatment. # (p < 0.05) and ## (p < 0.01) indicate significant differences from cells lacking Feijoa extract treatment.
Feijoa Peel and Pulp Extracts Reduce Reactive Oxygen Species Production in Red Blood Cells
The fluorescence probe DCF assay elucidated the protective role of Feijoa extracts against oxidative stress in RBC, as reported in Figures 3 and 4. ROS production increased nearly two-fold in Hg-treated RBC compared to the negative control. In contrast, co-incubation with 10, 20, 40, or 80 µg/mL of both peel and pulp acetonic extracts incrementally reduced ROS production in RBC. Similar to hemolysis, Feijoa peel extract prevented ROS production more potently than the pulp, and remarkably reduced fluorescence levels by approximately 50% at 10 µg/mL, to near control values at Statistical significance was calculated with one-way ANOVA followed by Tukey's test. ** (p < 0.05) indicates a significant difference from cells lacking HgCl 2 treatment. # (p < 0.05) and ## (p < 0.01) indicate significant differences from cells lacking Feijoa extract treatment.
Feijoa Peel and Pulp Extracts Reduce Reactive Oxygen Species Production in Red Blood Cells
The fluorescence probe DCF assay elucidated the protective role of Feijoa extracts against oxidative stress in RBC, as reported in Figures 3 and 4. ROS production increased nearly two-fold in Hg-treated RBC compared to the negative control. In contrast, co-incubation with 10, 20, 40, or 80 µg/mL of both peel and pulp acetonic extracts incrementally reduced ROS production in RBC. Similar to hemolysis, Feijoa peel extract prevented ROS production more potently than the pulp, and remarkably reduced fluorescence levels by approximately 50% at 10 µg/mL, to near control values at the highest extract concentration. At the same concentrations, Feijoa pulp extract also significantly reduced ROS production compared to the non-Feijoa protected RBC, reaching a maximum of about 50% at 80 µg/mL. the highest extract concentration. At the same concentrations, Feijoa pulp extract also significantly reduced ROS production compared to the non-Feijoa protected RBC, reaching a maximum of about 50% at 80 µg/mL. Statistical significance was calculated with one-way ANOVA followed by Tukey's test. ** (p < 0.05) indicates a significant difference from cells lacking HgCl 2 treatment. # (p < 0.05) and ## (p < 0.01) indicate significant differences from cells lacking Feijoa extract treatment. the highest extract concentration. At the same concentrations, Feijoa pulp extract also significantly reduced ROS production compared to the non-Feijoa protected RBC, reaching a maximum of about 50% at 80 µg/mL. Statistical significance was calculated with one-way ANOVA followed by Tukey's test. ** (p < 0.05) indicates a significant difference from cells lacking HgCl 2 treatment. # (p < 0.05) and ## (p < 0.01) indicate significant differences from cells lacking Feijoa extract treatment. Statistical significance was calculated with one-way ANOVA followed by Tukey's test. ** (p < 0.05) indicates a significant difference from cells lacking HgCl 2 treatment. # (p < 0.05) and ## (p < 0.01) indicate significant differences from cells lacking Feijoa extract treatment.
Peel and Pulp Extracts Prevent Hg-Induced Glutathione and Membrane Thiol Depletion in Red Blood Cells
GSH depletion is a key mechanism of Hg toxicity due to the weakening of the antioxidant defense system. We therefore evaluated the possible protective effect of Feijoa peel and pulp extracts on this specific, Hg-induced metabolic condition. As shown in Figure 5, 4 h treatments of RBC with 40 µM HgCl 2 reduce GSH levels by about 40%. Co-incubation with 20, 80, or 100 µg/mL of Feijoa peel extract prevented GSH depletion by about 20% with each concentration, such that GSH levels were unchanged from healthy control levels at the latter two peel extract concentrations. For the pulp extract, data indicate that 20 µg/mL had no effect on GSH levels, while significant protection was observed at 80 and 100 µg/mL, to a maximum of about 90% GSH levels compared to controls. RBC after incubation with HgCl2 ( Figure 6). Exposure to 40 µM HgCl2 reduced the level of membrane thiols by about 40%. This depletion was significantly counteracted by about 45% and 75% given co-incubation with 40 and 80 µg/mL of peel extract, respectively. Again, the pulp extract was less protective than the peel at the same concentrations, such that membrane thiol depletion was counteracted only at 80 µg/mL, by about 50%. Based on these results, it seemed appropriate to evaluate the efficacy of peel and pulp extracts in reducing the Hg-induced depletion of membrane thiols, using membranes obtained from intact RBC after incubation with HgCl 2 (Figure 6). Exposure to 40 µM HgCl 2 reduced the level of membrane thiols by about 40%. This depletion was significantly counteracted by about 45% and 75% given co-incubation with 40 and 80 µg/mL of peel extract, respectively. Again, the pulp extract was less protective than the peel at the same concentrations, such that membrane thiol depletion was counteracted only at 80 µg/mL, by about 50%.
Peel and Pulp Extracts of Feijoa Reduce Microvesicles Released from Red Blood Cells
To investigate the protective role of Feijoa extracts on erythrocyte morphological changes and MV formation known to be induced by Hg treatment [41,46], cells treated with HgCl2 and peel or pulp extracts, as described in the Materials and Methods section, were analyzed with confocal 9). Statistical significance was calculated with one-way ANOVA followed by Tukey's test. * (p < 0.05) and ** (p < 0.01) indicate significant differences from cells lacking Feijoa extract treatment.
Peel and Pulp Extracts of Feijoa Reduce Microvesicles Released from Red Blood Cells
To investigate the protective role of Feijoa extracts on erythrocyte morphological changes and MV formation known to be induced by Hg treatment [41,46], cells treated with HgCl 2 and peel or pulp extracts, as described in the Materials and Methods section, were analyzed with confocal microscopy. Hg treatment was associated with loss of the typical erythrocyte biconcave shape, as well as the formation of MV clearly discernible on cell membranes (not observable in the control) (Figure 7, Panel A). Cell treatment with Feijoa extracts completely restored the typical biconcave shape at 20 µg/mL and 80 µg/mL for peel and pulp, respectively ( Figure 7C-F).
Discussion
Mercury is not only highly toxic, but is an increasingly pervasive dietary heavy metal. As a matter of fact, concerns on the effect of Hg exposure on human health are not only limited to occupationally exposed workers, but also to the general population, mainly via contaminated food ingestion. Although in some European populations the overall Hg daily intake is below the tolerable amount [50,51], appreciable proportions of large fish populations are reported to contain levels of this heavy metal exceeding this amount, up to 2.22 mg/kg wet weight, including anglerfish (Lophius piscatorius) and black-bellied angler (Lophius budegassa) [52]. Discovering analogous means to simultaneously combat and protect against diet-based Hg toxicity is therefore crucial to public
Discussion
Mercury is not only highly toxic, but is an increasingly pervasive dietary heavy metal. As a matter of fact, concerns on the effect of Hg exposure on human health are not only limited to occupationally exposed workers, but also to the general population, mainly via contaminated food ingestion. Although in some European populations the overall Hg daily intake is below the tolerable amount [50,51], appreciable proportions of large fish populations are reported to contain levels of this heavy metal exceeding this amount, up to 2.22 mg/kg wet weight, including anglerfish (Lophius piscatorius) and black-bellied angler (Lophius budegassa) [52]. Discovering analogous means to simultaneously combat and protect against diet-based Hg toxicity is therefore crucial to public health. In this respect, phytochemicals able to counteract structural and metabolic alterations associated with heavy metal exposure are attractive for the reduction of their toxicity [53][54][55][56]. Data from our group indicate that hydroxytyrosol, an olive oil-derived phenolic antioxidant, has the potential to modulate the toxic effects exerted by Hg in human RBC [41][42][43].
To expand data on the potential role of nutrition in heavy metal toxicity, intact human RBC were exposed to 40 µM HgCl 2 , in line with our previous studies. Several markers of cellular toxicity were then evaluated to test the protective effect of Feijoa fruit extracts. According to data reported in similar experimental conditions, RBC treatment with 40 µM HgCl 2 for 4 h results in a doubling of ROS production, as indicated by DCF fluorescence [41]. Hg-induced ROS generation follows a significant decrease of GSH, which builds up a pro-oxidative microenvironment and renders cells more susceptible to ROS-mediated oxidative damage. A significant decrease in membrane thiols is also detectable in Hg-treated cells. The resulting hemolysis is significantly increased and measurable later, at 24 h.
Here we show the first evidence of Feijoa fruit extract protection against HgCl 2 -induced toxic effects in human RBC. The acetonic extracts of both the pulp and peel were able to counteract oxidative stress and cellular thiol decrease in Hg-treated RBC. The peel extract had a greater protective effect compared to the pulp, although to varying extents for the different markers analyzed, which is at least partially due to the greater proportion and diversity of polyphenols in the peel [3,57,58]. Interestingly, the protective effect of the peel from ROS production is only two-fold, compared to an eight-fold effect against overall cytotoxicity indicated by hemolysis. Whereas Hg sequesters and inactivates GSH by binding to sulfhydryl groups, polyphenols act on the resulting ROS by virtue of their hydrogen and electron transfer abilities. The presence of additional bioactive compounds with different activities in the peel (i.e., chelating properties) can also be hypothesized. Remarkably, as little as 10 µg/mL of Feijoa peel extract significantly affects all the tested markers, and 80 µg/mL completely prevents Hg-induced ROS production.
The data presented in this paper, although obtained from in vitro studies on human cells, also offer significant experimental evidence that Feijoa extracts prevent Hg-induced RBC shape alteration in RBC, which could be taken into account for future clinical investigations. In fact, although a particularly high GSH concentration may partially protect RBC from Hg's toxic effects, chronic exposure could affect RBC viability and induce morphological changes, also affecting cardiovascular disease. As mentioned before, Hg exposure enhances pro-coagulant activity of these cells, resulting in a contributing factor for Hg-related thrombotic disease [46]. In our previous studies, we raised the fascinating hypothesis that metabolic and shape modification of RBC may be regarded as a clinical biomarker, indicating increased cardiovascular risk in Hg-exposed individuals [41,42].
Our findings, in agreement with the literature data, strengthen the nutritional relevance of Feijoa bioactive compounds to the claimed health-promoting effects of this fruit. There is growing interest in utilizing Feijoa fruit for human consumption, due to its appetizing quality and its claimed health benefits. Feijoa fruit is an excellent source of vitamins and nonessential nutrients, as well as a variety of bioactive compounds endowed with significant antioxidant, antibacterial, and anti-inflammatory activities [10][11][12][13]. In this respect, there is a general agreement that the health-promoting effects of fruit and vegetable intake result from the combined properties and synergistic action of all bioactive constituents, including polyphenols [59,60]. These compounds can improve health due to their strong antioxidant activity, counteracting oxidative stress-induced cellular dysfunctions and modulating key mechanisms implicated in the development of oxidative stress-related human pathologies. Polyphenols are very useful in combating the deleterious effects of heavy metals. For example, Sobeh et al. [61] isolated and identified two compounds from the leaves of Syzygium samarangense (myricitrin and 3,5-di-O-methyl gossypetin), both showing antioxidant activities [62,63] and strongly reducing intracellular ROS accumulation and carbonyl content, while also protecting the intercellular GSH levels in keratinocytes (HaCaT) after exposure to sodium arsenite, one of the more toxic environmental heavy metals [61].
Feijoa has been proposed as an ideal candidate for nutraceutical strategies in the development of functional foods [64]. The data reported in this paper expands upon the known beneficial effects of Feijoa fruit, particularly related to chronic human exposure to heavy metal. In this respect, an interesting observation is that the very low active concentrations utilized in our study could be approached in vivo upon daily intake of Feijoa fruit. In this respect, some studies indicate that Feijoa fruit extracts are well tolerated in animal models. Karami et al. [65] demonstrated the hepatoprotective activity of methanolic extract of Feijoa fruit in a concentration range of 10-100 mg/kg, using the isolated rat liver perfusion system. The same group also investigated nephroprotective effects of leaf extracts (10-40 mg/kg) on renal injury induced by acute doses of ecstasy (MDMA) in mice [66]. Moreover, in a recent study, Feijoa leaf extract was shown to be devoid of toxicity in rats up to 2 g/Kg, [11]. Finally, we have confirmed by MTT test, on human leucocytes as well (data not shown), that treatment for 24 h with 5, 50, and 500 µg/mL of acetonic extracts of F. sellowiana did not induce significant cytotoxic effects, as already demonstrated on either the Caco-2 or HT-29 cell lines [16].
The food industry is increasingly interested in the utilization of non-edible parts of fruits. Phytochemicals are proposed for designing foods with added functional value, aiming to beneficially affect target functions in the body and reduce the risk of diseases. These compounds are present in large quantities in waste products from the agri-food supply chain, especially peels and seeds. Our data, showing a greater protective effect from the Feijoa peel on Hg cytotoxicity than from the pulp, corroborates this rationale. Recovering and using such a waste product, normally destined to magnify industrial waste production, would give new life to the less noble part of the fruit. This is further in line with recent studies that propose the potential utilization of Feijoa fruit peel for added processing and functional value. As demonstrated by Sun-Waterhouse et al. [64], the extracts produced from Feijoa waste, such as the peel, retain high pectin content, which is advantageous for food applications. Moreover, the possibility to utilize Feijoa peel-containing food packaging film for the inhibition of foodborne bacteria was recently demonstrated [67].
In conclusion, the novel beneficial properties of Feijoa reported in this paper, regarding its efficacy to reduce heavy metal toxicity in human RBC, provide biochemical bases for the use of Feijoa-based functional foods or pharmacological preparations in preventing and combating mercury-related illnesses. | 5,829.6 | 2019-07-01T00:00:00.000 | [
"Biology"
] |
Sustainable Development of Livestock and Meat Production in Republic of Benin : Strategies and Perspectives
| This paper aims to update information for a better understanding of the functioning of the sector of animal and meat production in Benin’s socio-economic context. It highlights (i) structure of livestock population and production, (ii) competitiveness of meat sector, (iii) identifies constraints and possible prospective solution to increase meat production in Republic of Benin. The traditional animal production systems remain largely widespread. However, industrial and modern livestock farming systems for all species are developing. Cattle (57%), chicken (19%), small ruminants (13%) and pigs (7%) are the main meat producers in the country. Non-conventional species such as rabbit (3%) and grass cutters (1%) contribute also to the national meat production. Despite religious prohibitions, pork consumption increased during the last years notably in southern Benin. The households with a higher monthly income spend more money to purchase meat than poor households. Taste, texture, price, and juiciness are the main criteria of choice. Also, consumers prefer the meat of local breeds to that of exotic breeds. Thus local species and breeds have a great role in the development of this sector. Although policies have been implemented to boost the national meat production, the sector is still undeveloped. That is why, the implementation of new approaches and practices including improvement of animals’ genetic resources, housing, health care, and feeding should be developed to intensify production.
IntroductIon
T he Republic of Benin is a West African country, bordered by Togo to the west, Nigeria to the east, and Burkina Faso and Niger to the north. The territory of Benin covers an area of 114,763 square kilometres with estimated a population to be approximately 10.08 million in 2013 (NISEA, 2015). Land resources consists of forest (40%), agricultural land (31.3%: arable land 22.9%; permanent crops 3.5%; permanent pasture 4.9%) and other (8.7%) (CIA, 2017). With 32.7%, the agricultural sector is the main contributor to Benin's Gross Domestic Product (GDP) and employed 75% working population of the local population. The subsector of animal production represents nearly 13% of GDP (Sodjinou et al., 2007;MALF, 2013).
In the Republic of Benin, the main product from livestock production is meat. The livestock population comprises conventional species including poultry, pigs, small ruminants, and cattle which are the major source of protein for the local population. It also comprises non conventional species made of grass cutters (Cane rats), snails and rabbits. For most species, animals are reared in traditional production system. However, improved (semi-intensive and intensive)systems are being developed notably in avian sector. Meat production does not meet the expressed needs of consumers and leadsto an increasing meat importation every year (Youssao et al., 2008). According to the National Office for Livestock (NOL) (Direction de l' élevage in French) (2005), the demographic growth in the Republic of Benin leads to a steadily increase of demands in animal proteins notably meat. The Food and Agriculture Organization (FAO) (2005) and NOL (2014) reported a level of meat consumption estimated at 12 kilograms per capita per year. This level is largely lower than that recommend by FAO (21 Kg/per capita/years). Furthermore, the increasing socio-economic conditions of the consumers in the developing countries can be accessed from changes in their consumption pattern. Therefore information about consumers' meat preference is crucial in developing and implementing appropriate livestock improvement strategies (Ogbeide, 2015). This paper highlights briefly the present situation of livestock and meat production in the Republic of Benin.
MEtHodology
We wrote this paper using desk research methodology to compile information from scientific papers various studies reports, strategy documents obtained from the Ministry of Agriculture, Livestock and Fishery of Republic of Benin or from the website of Food and Agriculture Organization of the United Nations (FAO) and statistical data. A total of eight articles searched online using the keywords including Benin, meat, animal, systems of production, and constraints were used when writing this paper. In addition, some website including: Food and Agriculture Organization of the United Nations-Division of Statistics (FA-OSTAT), memoire online, FAOSTAT when collecting information and data. We used statistical data related to cattle, small ruminants, pigs and poultry (Chickens broilers and layers) population and meat production from the annual report and statistics of National Office for Livestock (NOL) and Ministry of Agriculture, Livestock and Fishery of Republic of Benin (MALF) (2007)(2008)(2009)(2010)(2011)(2012)(2013) and the FAOSTAT (1990FAOSTAT ( -2006FAOSTAT ( and 2014. Using these collected data, we performed descriptive statistical (annual growth average (AAGR), Livestock Unit (LU), charts, and percentage) analysis using Microsoft Excel 2007. We calculated AAGR using geometric formula (FAO, 2011) and LU as described by FAO (FAO, 2005).
Where: X t is the final value of population;X 0 is the starting value of population; (t-0) is the number of years.
rESultS And dIScuSSIon livesTock PoPulATion
The animals' population trends from 1990 to 2014 are given in the Table 1. The population of cattle estimated at over 2.2 million heads in 2014 is made of 31% trypanotolerant cattle (Borgou, Somba and Lagunaire), 7.7% zebus (M'bororo, Goudali and White Fulani) and 61.3% crossed breeds (Alkoiret, 2011;Hervé, 2017). These local breeds are very well adapted to adverse climate and breeding conditions. They have the ability to walk long distances in search of feed and high resistance to diseases and parasites. However, studies showed that the low genetic potential, resulting in low milk and meat production, required the implementation of appropriate breeding programs (Dehoux and Hounsou-Ve, 1993;Adjou Moumouni, 2006;Alkoiret et al., 2011). Thus, in order to increase national milk production and to meet increasing milk needs of the population, Brazilian Girolando cattle have introduced in November 2004. This exotic breed producing 7.22 ± 0.15 Kg milk per day is kept in some state and private farms (Doko et al., 2012). The number of small ruminants also increased from 1.58 million heads in 1990 to 2.63 million heads by 2014 with an annual growth rate of 2.16% in 2014. The population of both of two species (goats and sheep) is largely constituted by the West African Dwarf goat and sheep breed (Also known as the Guinean or Djallonké breed) (Aregheore, 2009; Adote et al., 2011;Monkotan, 2011). West African dwarf goat and sheep are kept essentially for meat production. They are well adapted to conditions in humid and sub-humid zones characterized by strong presence of tsetse flies. They are small size animals with a low meat yield and very low lactogenic productive potential (Gbangboché et al., 2002;Monkotan, 2011). Guinean sheep live weight is rarely more than 30 kg with a yield of about 48% but has a high quality of meat (Aregheore, 2009). The second breed raised in the Republic of Benin is Sahelian breed. It is mostly found in the extreme north of Benin. The Sahe-
Journal of Animal Health and Production
March 2018 | Volume 6 | Issue 1 | Page 37 lian sheep are large and heavy (80 kg) animals. They are meat-purpose animals (40 to 50% of yield) (Aregheore, 2009). Additionally, Alpine goat breeds have been imported in order to diversify milk production, which still comes from cows only. They are currently being kept under an acclimatization research program in some private farm (Monkotan, 2011;Vissoh et al., 2015).
Pig population which was 462 thousand heads in 1993 drastically decreased nearly 255 thousand heads in 2000, and increased more than 431 thousand animals. This decline was due to African swine fever outbreak in 1997. This epizootic decimated more than 70% of national pigs' population causing huge economic losses (Ayssiwede, 2004).
Most of the pigs are predominantly local breeds. Despite their low growth and small size, local pigs have a greater ability to withstand disease, hard local breeding, and environmental conditions. Apart from these native breeds, exotic pigs (including Large White, the Landrace and their crossbreeds) are also found. They are large size animals raised in semi-modern or modern systems (Youssao et al., 2008).
Poultry farming is a widespread production activity throughout Benin. It is subdivided into two types: traditional farming in which local breed commonly known as poulet bicyclette (Sahouè, Fulani, Koungbo, etc) are reared. This system largely preponderant aims self-consumption of families and marketing. It is characterized by the use of low inputs. There is also modern poultry farming which exploits improved or imported poultry breeds such as ISA Brown, Warren, Hyline, Lohman, Rhode Island Red and Plymouth Rock (Sodjinou, 2009;FAO, 2015). In 2013 of the estimated 18.2 million heads, only 1.2 million heads are improved or exotic breeds (NOL, 2014).
Most small and large-scale commercial farms are located in southern Benin. As we move northward, poultry farms become scarcer, with priority given to local poultry and to cattle and goat farming (Ogbeide, 2011). The grasslands of the northern Benin are the main hub of ruminants' production in Benin. Indeed, the large population of cattle and small ruminants is kept in Borgou-Alibori with nearly 63% and 29 % of the population respectively. As for chicken and pig, they are mostly found in Zou-Colline and Ouémé-Plateau with 46% and 27% respectively.
This geographical distribution of population may be explained by the variation of climatic factors from one region to another. These climatic conditions influence the availability of food resources and the pressures of tsetse flies vectors of trypanosomiasis (Alkoiret et al., 2011). It can also due to dietary behaviours of consumers, targeted markets, social and religious factors (Islam) (Ayssiwede, 2004).
AnimAl ProDucTion sysTems
Generally, animal production systems are numerous but the main production systems are traditional. This diversity is mainly due to range of agro-ecological zones in Benin, social and ethnic groups and the technical level of keeping. Animals are mostly kept under traditional production system. These systems include transhumant-pastoral, sedentary (free roaming), urban and peri-urban and sometime semi-extensive breeding system.
TrAnshumAnT-PAsTorAl sysTem
It is the most common system of production in northern and central-northern of Benin where cattle are the main species of breeding. Small ruminant flocks are secondarily annexed to that of cattle. This system is mostly practiced by the Fulani and Gando pastoralists whose livestock keeping is the principal source of income (NOL, 2005). This system is a seasonal and cyclical movement (sometimes cross-border) of keepers and herds according to the rainfall regime (Adjou Moumouni, 2006). It is characterized by the absence of production target, limited food resource, low productivity and high losses due to accidents, diseases, and theft. This system is actually being turned into a mixed farming system in which producers combine animal production with agriculture notably cultivation of cotton and some food crops (NOL, 2005;Ajala et al., 2008;Alkoiret et al., 2011). It generates conflicts between crop farmers and breeders due to non-respect of transhumance corridors causing crop damage.
seDenTAry sysTem
This system is a form of natural pasture, browse, crop residues, and kitchen wastes exploitation. In this system, animals are kept all year round on a fixed area (village) where livestock and crop production are mixed (Adjou Moumouni, 2006). They freely graze on the common village pasture and roam (free roaming) about scavenging for food over the day. At night, they return to their sheds. This system includes crops farmers and non-rural people who entrust their animals to sedentary keepers who and share the offspring. It is the predominant system in central, southern regions and West-Atacora (NOL, 2005). Additionally, family labours play actively an important role in livestock management and animals may not be offered sufficient feed, clean water, and healthcare.
urbAn AnD Peri-urbAn sysTem
This system is more found in the urban peri-urban areas, where animals are confined in a yard or barn for fattening. The main sources of feed for animals are agro-industrial by-products, kitchen wastes from household food preparation and forages along roadsides and undeveloped plots in the cities (Baah et al., 2012).
Journal of Animal Health and Production
March 2018 | Volume 6 | Issue 1 | Page 38
nATionAl meAT ProDucTion
Although the increase of animals numbers, research and strategic policy efforts of the government to improve national production, meat production notably remains underdeveloped. Meat is mostly produced from conventional species including bovine, small ruminants, pigs, and poultry. However, the rearing of non-conventional species (snails, grass cutters, and rabbits) contributes also to national meat production. As shown in Figure 1, the indigenous meat production increased from 36,290 tonnes in 1993 to 64,968 tonnes in 2013 with a yearly average of 46,706 tonnes. Cattle are the main contributors to national meat production with 57%. It is followed by chicken (19%), small ruminants (13%) and pigs (7%). As for non-conventional species, rabbit and grass cutters meat represented less than 3% and 1% respectively of total production. In 2013, this production represented 2.25% of West African's indigenous meat production that was over 2.78 million tonnes. Among West African countries, the largest amount of indigenous meat is produced in Nigeria (1.22 million tonnes), Mali (356 thousand tonnes) and Niger (252 thousand tonnes) accounting for 44.11%, 12.83% and 9.06% respectively (FAOSTAT, 2017b).
meAT imPorTATion AnD exPorTATion
Republic of Benin has considerable difficulties meeting the needs of its populations in animal protein, especially meat. Thus, to meet this deficit the country imports from neighbouring or European countries animals and frozen meat. Conforming to the statistic of the National Office for Livestock in 2013, more than 177 thousand tonnes of frozen meat has been imported. However, only about 10% of imported frozen meat is consumed by the local population (NOL, 2014). The most important shares are exported neighbouring countries notably Nigeria where high demand is noticed. This may be justified by the local consumers' preferences toward the meat locally produced. Small ruminants are mainly imported from Nigeria and Niger, while Niger and Burkina Faso provide essentially cattle. As for frozen meat imported, it mostly concerns poultry meat and poultry offal. Among countries exporting poultry meat and poultry offal to Republic of Benin, Brazil (15.21%), Spain (13.34%), Poland (12.98%), France (11.78%), and UK (10.31%) are the most important. In 2013, of the total imported over 38.23 thousand tonnes (including meat from imported animals and frozen meat) have locally been consumed. Cattle, small ruminants, pigs with over 17.54, 2.22 and 0.17 thousand tonnes respectively were most important meat from imported animals. While frozen poultry meat represents the most shares of frozen imported meat with nearly 17.65 thousand tonnes accounting for 96.39 % of total frozen meat. The frozen meat of small ruminants was estimated at 0.46 tonne and accounted for a negligible share 0.003 %. Cattle, pigs, and rabbit represent 2.93%, 0.32%, 0.34% respectively of total frozen imported meat (NOL, 2014).
AnimAls mArkeTing, meAT Processing AnD consumPTion
As shown in Figure 2, live animals are marketed on the local market directly by owners or their family members or indirectly by retailers. Studies showed that the main criteria that determine animal prices in these traditional markets are breed, the time of the year, the region, weight, sex, and age (Sodjinou et al., 2007). Indeed, as one draws closer to urban areas, animals' prices rise. Additionally, low prices are mainly observed at the beginning the school year, where the keepers need cash to pay tuition fees of their children and to start crops production activities (Sodjinou et al., 2007). Also, an increase of sheep price is observed during celebrations period such as the Muslim feast also known as Eid-el-Kabir during which, a high demand is noticed. Animals sold are transported and slaughtered by an individual (unlicensed slaughtering) or in the state, municipal, breeder association or private slaughterhouses or slaughtering floors. Then, the cuts of fresh meat are collected and sold by retailer butchers to consumers on the local markets. On another hand, the fresh meat is processed and
Journal of Animal Health and Production
March 2018 | Volume 6 | Issue 1 | Page 39 Approaches of Solution • Low productivity of local breeds that constitute the large part of the animal population: although local breeds are very well adapted to environmental and breeding conditions, they have the low genetic potentials and productivity.
• More research should be done to promote the implementation of an effective breeding and genetic improvement program of indigenous breeds through the application of assisted reproductive technologies such as artificial insemination, semen cryopreservation etc.
• The lack of fund for research in livestock production and development, and support the cooperation between research institutions (University, research centres etc) and farmers association to promote a great genetic evolution of the herds and animal development.
• The government should provide more funds to support local farmers, researchers and genetic improvement centres which should provide high genetic value reproducer animals to breeders in order to improve the animal productivity in real farm area and ensure a better monitoring of genetic progression.
• The inadequacies of animal identification, the system for collecting statistical data and lack of updated official data on livestock population across the country.
• An effective animals' identification system using ear-tags and registration (computer or internet based systems in regional livestock directorate) should be set up to ensure an easier management and availability of updated statistical data.
• Poor animal health monitoring due to the insufficient of veterinary services. This is expressed through a high incidence of infectious, parasitic and viral diseases in all species with enormous financial losses for keepers notably in traditional systems. .
• The government should implement an appropriate animal health monitoring (animal health inspection service) with a view to strengthening diseases and epizootics eradication measures.
• The absence of long-term development policies and strategies of livestock sector due to lack of governmental funds and awareness for sustainable development of animals farming activities. Difficulty of famers to access banking loans which are at high interest rate.
• The government should implement better animal development assistance policies by providing technical and financial assistance to farmers to improve the production systems. It should reduce the meat importation by increasing the custom taxes and easing the access of farmers to the banking loans at low and effective interest rate.
• The lack of training for animal breeders in term of modern technologies of animal reproduction and farm management.
• The regional livestock directorates should regularly set up different training session in order to educate and update the knowledge of breeders about improved techniques of animal production systems and flock management strategies.
sold by restaurants or some individuals as grilled meat directly consumed.
The households with a high monthly income spend more money than poor households to purchase of local chicken (or meat). Taste, texture, price, and juiciness are the main criteria of choice (Sodjinou et al., 2007). Governmental authorities through Ministry of Agriculture, Livestock and Fisheries and the Ministry of Trade and local municipal administrations regulated and controlled the livestock and meat price on the markets. Also, The Ministry of Agriculture, Livestock, and Fishery in collaboration with Ministry Public Health ensure hygienic quality control of meat and population health.
consTrAinTs of livesTock ProDucTion AnD APProAch of soluTion To boosT The AnimAl ProDucTion AnD meAT secTor
Although meat sector is the most exploited way of livestock resource, meat production, and marketing showed some difficulties due to livestock management, lack of structural and institutional organization of the subsector. The major problems associated with livestock production and some possible solutions to increase the meat production are presented in the Table 2. concluSIon Livestock production plays an important economic and social role. In Republic of Benin, livestock population and meat production increased over last decades. However, the national needs in meat products did not meet due to many limiting factors including the low genetic value of animals, production systems and lack of governmental policies. This leads to an importation of animals and animal products (cuts of meat, poultry offal) from others countries. Given that the human population and the needs in animal protein are increasing, efforts should be made to improve livestock genetic improvement, production systems in order to significantly increase national production. Also, the quality of meat consumed has to be more controlled to protect consumers against health problems.
AcKnowlEdgEMEntS
The authors are thankful to the National Office for Live- | 4,657.8 | 2018-01-01T00:00:00.000 | [
"Agricultural and Food Sciences",
"Environmental Science",
"Economics"
] |
Smart Contract Design Pattern for Processing Logically Coherent Transaction Types
: Recent research shows that the source code of smart contracts is often cloned. The processing of related types of transactions in blockchain networks results in the implementation of many similar smart contracts. The rules verifying transactions are therefore duplicated many times. The article introduces the AdapT v2.0 smart contract design pattern. The design pattern employs a distinct configuration for each transaction type, and verification rule objects are shared among configurations. The redundancy of logical conditions was eliminated at two levels. Firstly, it is possible to combine similar smart contracts into one. Secondly, a configuration in a smart contract reuses verification rule objects at runtime. As a result, only one object is instantiated for each verification rule. It allows for the effective use of operating memory by the smart contract. The article presents the implementation of the pattern using object-oriented and functional programming mechanisms. Applying the pattern ensures the self-adaptability of a smart contract to any number of transaction types. The performance tests were carried out for various numbers of verification rules in a smart contract and a different number of checked transactions. The obtained evaluation time of 10,000,000 transactions is less than 0.25 s.
Introduction
Smart contracts are software that controls the execution of transactions in blockchain networks.An overview of smart contracts was presented by Zheng et al. [1] with a discussion on the challenges and recent technical advances.They reveal that smart contracts are written mostly in the following programming languages: Solidity, Go, Java, and Kotlin.They also indicate problems that plague current blockchain platforms, i.e., re-entrance, block randomness, and overcharging.They emphasize that it stems from the under-optimization of smart contract source code, in which various anti-patterns may be found (e.g., dead code or expensive operations in loops consisting of repeated computations).Blockchain technology is increasingly used in a wide range of applications.Hence, the topic of designing, programming, and testing smart contracts is becoming more and more important.Lately, Wu et al. [2] reviewed the progress that has been made in smart contracts.They confirmed the essence of the design process, and showed the smart contract life cycle including the phases: contract generation, contract release, and contract execution.The authors pointed out the main problems in the development of smart contracts in the areas of performance, privacy, and security.As for efficiency, they underlined that the problem lies in low contract execution efficiency.Moreover, Kannengießer et al. [3] identified challenges in smart contract design, proposed solutions, and recommended software design patterns.Their recommendations sometimes refer to existing software design patterns previously proposed by Gamma et al. [4], e.g., Proxy Pattern and Façade Pattern.However, they also show smart contract-specific patterns.Researchers and practitioners propose using the Oracle Pattern whenever external data are required by a smart contract.Additionally, the known threat in smart contracts is a reentrancy attack.The authors propose using Mutex Pattern or Checks-Effects-Interactions Pattern to eliminate that smart contract vulnerability.In the area of efficiency, Six et al. [5] pointed out patterns that enhance execution, storage, and redundancy aspects, e.g., Incentive Execution, Limit Storage, and Avoid Redundant Operations.In addition, researchers work on different ways to execute smart contracts on blockchain networks.Recently, Liu et al. [6] proposed parallel processing of transactions that increased the level of throughput.Design patterns are constantly evolving.Gupta et al. [7] proposed Proxy Smart Contracts that serve as intermediaries in the execution of the actual smart contracts.That pattern is used for the communication of on-chain smart contracts with off-chain services.Another software design pattern that has been employed for blockchain smart contracts is the Delegation Pattern.Kim et al. [8] have applied it to the construction of updatable smart contracts to provide federated authentication schemes.Additionally, Proxy and Delegation patterns are typically used in the upgrading mechanism of Ethereum smart contracts [9].A different manner of standardizing the design of smart contracts is applying templates.Chu et al. [10] review works that not only identify vulnerabilities, but also offer mechanisms and tools to eliminate them.They show the currently developed mechanisms used to repair vulnerabilities of smart contracts operating in off-chain and on-chain modes.
However, the currently developed design patterns for smart contracts do not cover the subject of their reconfiguration, in particular the possibility of adapting to various types of processed transactions.A smart contract can perform the same operation for logically consistent but different types of transactions.That may include, e.g., on-chain and off-chain transactions, authentication and authorization in the system for different roles, in-community and cross-community energy transfers, and domestic and foreign contract management for the cross-border labor market.Hence, a need to define a smart contract design pattern that would enable operations on various types of transactions.Górski presented the design pattern for reconfigurable smart contracts in [11].However, it only allowed operation on two types of transactions.In addition, the reconfiguration of the smart contract entailed the need to create new objects of verification rule classes.As a result, it engaged the garbage collector every time the transaction type changed.Additionally, the abstract layer of the pattern was not entirely independent of the specific smart contract.
Figure 1 illustrates smart contract configurations for transaction types, which share a set of unique verification rule objects.The following symbols are used in the figure: obj i means the object of the i-th verification rule and re f i means a reference to the obj i .Verification rule objects are shared within one smart contract, which allows checking logically related transactions.The main contributions of this paper are listed as follows: • The AdapT v2.0 design pattern that allows processing any number of transaction types; • A redesign of the abstract layer of the pattern; • An implementation of the abstract layer in Java language, which employs objectoriented and functional programming.The implementation of the abstract layer is independent of the specific smart contract; • Implementation of the concrete layer of the pattern in Java language on the example of a smart contract for energy transmission in prosumer communities; • Reuse and redundancy metrics with analysis of the pattern; • Performance tests of the pattern for the number of transactions ranging from 100,000 to 10,000,000.
The remaining part of this paper has the following structure.Section 2 presents related studies on the topic tackled in the paper.Section 3 introduces the design of the AdapT pattern.The section also contains the implementation of the pattern in Java.In Section 4, an analysis of the re-use and redundancy was enclosed, whereas Section 5 depicts the results of performance tests.Section 6 contains a discussion on the pros and cons of the pattern.Section 7 summarizes the work performed and lists the planned future tasks.
Related Work
The topic of smart contract development is a very fast-growing branch of software engineering.Despite the fact that significant progress has been made in the area of smart contract design, many challenges remain.The interview among professionals performed by Zou et al. [12] reveals the following obstacles that developers must face: lack of design methods of secure smart contracts, basic existing development tools, limitations of programming languages, performance constraints, and limited online resources.Vacca et al. [13] investigated tools developed for smart contracts.They highlight that the majority of the tools target the Ethereum framework, e.g., SolMet, GasChecker, and SmartEmbed.Only a few tools were written for Hyperledger Fabric, e.g., Zeus and Blockbench.The tools gathered in the review help mainly detect security vulnerabilities in smart contracts written in Solidity language.Two of them have functionality close to software design issues: SmartEmbed identifies code clones and Gasper locates gas-costly programming anti-patterns.Gas waste in smart contract loops is considered by Li et al. [14] focusing on applying machine learning to detect the Expensive Operation anti-pattern.However, none of them verifies the source code of smart contracts on compliance with design patterns.Two design patterns were introduced by Mandarino et al. [15].The first is an architecture that reduces gas consumption in case of updating the source code of a smart contract.The second shows smart storing of data in the form of packing bits representing boolean values.The pattern for communication of on-chain smart contracts with off-chain services was proposed by Liu et al. [16].They proposed a data carrier architecture that consists of three components: Mission Manager, Task Publisher, and Worker.Those components interact with smart contracts and off-chain data sources.
Researchers also propose templates as an alternative way to unify smart contract design.For example, Jin et al. [17] developed the Aroc tool, which generates a patch contract containing security rules based on the fixed template and deploys it to the blockchain.Templates are also the subject of other research.Furthermore, Gec et al. [18] propose a support system that recommends and provides smart contract source code templates suitable for a fog architecture, whereas Mao et al. [19] operate at a higher level of abstraction.They provide a set of specialized templates of basic functions for users to design smart contracts visually by the user interface.
An interesting area where design patterns may appear is model-driven engineering, because generating smart contract source codes requires some template or design pattern.In this context, Bodorik et al. [20] show the transformation that generates the source code of smart contracts from the description of business models presented in Business Process Model and Notation (BPMN).Similarly, Shen et al. [21] also use BPMN to visualize the process to be understood by the stakeholders from various domains.They introduced a smart contract generator for developing multi-party interaction scenarios.In contrast, Jurgelaitis et al. [22] step lower at the level of system modeling and show the mechanism of generating smart contracts for the Solidity language from Unified Modeling Language (UML) state diagrams.
Currently conducted research works concern the application of blockchain technology and smart contracts in various domains.Solutions for the medical sector are among the most intensively researched (Yang et al. [23]).Additionally, smart contracts are increasingly used in the energy sector, especially in distributed renewable energy systems (Honari et al. [24]).From its beginning, blockchain has been used in the financial sector (Wang et al. [25]).On the other hand, smart contracts are increasingly used in public services such as online voting systems (Saim et al. [26]).Logistics, in particular supply chain management, is also a constantly developed area of application (Natanelov et al. [27]).Interesting research works can be expected in the area of smart contract unification.There is a place for research on both domain-specific and domain-independent design patterns.Anyway, the first papers on these issues have already been published.Wohrer et al. [28] show how to move from domain-specific language to smart contract code.On the other hand, Capocasale et al. [29] discuss the issue of standardization of smart contracts in a broader context, independent of the field of application.Various formal verification methods are used in efforts to improve smart contracts.Concerning recent work in this area, Nam et al. [30] propose to analyze Solidity smart contracts using Alternating-time Temporal Logic model checking, whereas Almakhour et al. [31] deal with formal verification of composite smart contracts that require other smart contracts to be executed.They employ the finite state machine models to verify Solidity smart contracts.In contrast, Pasqua et al. [32] introduce the method that analyzes Ethereum bytecode and extract precise Control-Flow Graphs.
From the point of view of software engineering, it is important to be able to reuse already written source code.In this regard, Pierro et al. [33] propose to organize a repository of Ethereum smart contracts.Smart contract source code reuse is also discussed by Chen et al. [34].Their research reveals that about 26% of smart contract source code blocks are reused in 146,452 analyzed open-source Ethereum projects, whereas Khan et al. [35] introduce the topic of source code cloning of Ethereum smart contracts.In this aspect, Górski [11] presents a design pattern that enables the reuse of validation rules used in a smart contract.The usage of classes to define verification rules enables the reuse of rules within the same contract and between different contracts.It should be also emphasized that the self-adaptation characteristic is not sufficiently researched in the context of smart contracts.However, scientific work is emerging in this area.Singh et al. [36] propose a self-adaptive security approach for smart contracts based on Service Level Agreements to provide countermeasures to attacks.Looking even more broadly from the point of view of software architecture, smart contracts are considered an important type of IT system function.The 1 + 5 architectural view model for cooperating systems proposed by Górski in 2012 already included the Contracts view [37].However, it was only at the EUROCAST conference in February 2019 that the same author presented the context of using the Contracts view for smart contracts [38,39].In this model, the key is to obtain business justification for the functions considered in the IT system.The business process is modeled in the Integrated Processes view.A similar aspect was underlined at the OTM Conferences in October 2019 by Bagozi et al. [40].They showed the design of smart contracts identified from business process models of collaborating organizations.In their work, they also used a two-level model of smart contract design: abstract and concrete.This confirms the validity of the adopted architecture for the AdapT pattern.
The currently proposed AdapT v2.0 design pattern takes source code reuse to an even higher level.Validation rule objects are shared at runtime.Thanks to this, their redundancy was eliminated.In consequence, the efficiency of memory usage has been raised.The pattern is also now adapted to support any number of transaction types.As a result, the pattern has also been made more flexible as far as self-adaptation is concerned.
The Pattern Design and Implementation
Further considerations require clarification of the terms used in the paper, i.e., the verification rule, verification rule object list, smart contract configuration, and evaluation expression.The author has proposed the following definition of a verification rule (Definition 1).
Definition 1 (Verification rule).
A single logical condition imposed on a smart contract.
Smart contracts may employ numerous verification rules.Moreover, the author has introduced the following definition of a smart contract verification rule object list (Definition 2).Definition 2 (Verification rule object list).An ordered collection of all non-recurring verification rule objects for all verification rules used in the smart contract.
Where a smart contract supports multiple transaction types, verification rule configurations apply.Therefore, the author has put forward the following definition of a verification rule configuration (Definition 3).
Definition 3 (Verification rule configuration
).An ordered list of non-repeating verification rule references that point to verification rule objects appropriate to the transaction type.
The notions of verification rule configuration and configuration will be used interchangeably hereafter.All verification rules that constitute a configuration must be met for the transaction to be executed.Verification rules in the configuration are used by an evaluation expression.The author has formulated the following definition for the notion of the evaluation expression (Definition 4).
Definition 4 (Evaluation expression). A logical expression containing verification rules and logical operators that return a single boolean value.
The pattern was constructed in division into two layers: Abstract and Concrete.The split was used to introduce an abstraction layer, common to all smart contracts designed according to this scheme.The elements of the layer are independent of the implementation of a specific smart contract and are reused in each of them.
Abstract Layer
The Abstract layer consists of two abstract classes: AbstractTransaction and AbstractS-martContract.The Abstract layer consists of two abstract classes: AbstractTransaction and AbstractS-martContract.The AbstractTransaction class serves as a parent class for all specific transaction classes handled by the specific smart contract.All specific transaction classes must inherit from that abstract class, whereas the AbstractSmartContract class sets a template for all specific smart contract classes.The class declares a list of verification rule objects (rulesList variable).That list employs the Predicate functional interface which in turn operates on reference type AbstractSmartContract.The class also declares a list of verification rule configurations for various types of transactions (configurations variable).Such structure will enable two characteristics: handling various reference types inheriting from Abstract Transaction, and the use of lambda expressions for deferred execution.In addition, the Ab-stractSmartContract class declares the checkSC() method, which is used to verify the smart contract.It has one input parameter, which is a reference to the transaction object to be verified.This supports the later implementation of pure type, where the result hinges only on the input data.Depending on the verification result, this method returns a true or false logical value.The combination of inheritance from object-oriented programming and lambda expressions from functional programming will allow for processing different types of transactions in this one method.Since they are abstract classes, none of them can be instantiated.
Both classes from the Abstract layer of the pattern were implemented in Java language.The source code of the AbstractSmartContract class is shown in Listing 1.
Listing 1.The source code of the AbstractSmartContract class.package adapT ; import j a v a .u t i l .A r r a y L i s t ; import j a v a .u t i l .L i s t ; import j a v a .u t i l .f u n c t i o n .P r e d i c a t e ; Only standard classes and interfaces available in Java language are used in the implementation.Additionally, the code is written to be domain-independent, whereas the Concrete layer of the pattern uses classes implementation-specific for the concrete smart contract.
Concrete Layer
Smart contracts are widely used in energy applications.Their usages were recently reviewed by Vionis and Kotsilieris [41].The Concrete layer of the pattern uses the example of energy transfer between different stakeholders in the distributed renewable energy system.In such systems, energy can be exchanged between prosumers within the same community and between prosumers in different communities.Additionally, electricity can be sent to the power grid.
The schema Figure 3 shows a UML Use case diagram for the SendEnergy use case.The Use case diagram employs stereotype <<IntegratedSystem>> from the UML Profile for Messaging Patterns defined by Górski [42].The stereotype denotes actors representing applications external to the prosumer one.One additional class the Transaction has been defined to gather common attributes of transactions.The class counteracts the redundancy of attributes.The Transaction class is also abstract.One can only instantiate objects of the three specific transaction classes.The class that inherits from the AbstractSmartContract abstract class is responsible for handling various transaction types.In the presented example the ExchangeEnergyContract class inherits from that abstract class.In the constructor of the smart contract class, both the verification rule list and configurations for transaction types are initiated.
The constructor source code of the ExchangeEnergyContract class presents Listing 3.
Listing 3. The source code of the ExchangeEnergyContract class constructor.
p u b l i c ExchangeEnergyContract ( ) { // v e r i f i c a t i o n r u l e s r u l e s L i s t .add ( t −> ( ( T r a n s a c t i o n ) t ) .getSourceID ( ) ! = ( ( T r a n s a c t i o n ) t ) .g e t T a r g e t I D ( ) ) ; r u l e s L i s t .add ( t −> ( ( T r a n s a c t i o n ) t ) .g e t Q u a n t i t y ( ) > 0 ) ; r u l e s L i s t .add ( t −> ( ( T r a n s a c t i o n ) t ) .g e t S o u r c e S u r p l u s ( ) >= ( ( T r a n s a c t i o n ) t ) .g e t Q u a n t i t y ( ) ) ; r u l e s L i s t .
add ( t −> ( ( T r a n s a c t i o n C r o s s ) t ) . getSourceCommunityID ( ) ! = ( ( T r a n s a c t i o n C r o s s ) t ) . getTargetCommunityID ( ) ) ; r u l e s L i s t . add ( t −> ( ( T r a n s a c t i o n ) t ) . getTargetNeed ( ) >= ( ( T r a n s a c t i o n ) t ) . g e t Q u a n t i t y ( ) ) ; r u l e s L i s t . add ( t −> ( ( T r a n s a c t i o n G r i d ) t ) . g e t T a r g e t I D ( ) == ( ( T r a n s a c t i o n G r i d ) t
) .getEnergySubnetID ( ) ) ; // c o n f i g u r a t i o n s f o r ( i n t i = 1 ; i <= 2 ; i ++) c o n f i g u r a t i o n s .add ( new A r r a y L i s t < >() ) ; // c o n f i g u r e r u l e s f o r T r a n s a c t i o n I n c o n f i g u r a t i o n s .g e t ( 0 ) .add ( r u l e s L i s t .g e t ( 0 ) ) ; In the constructor of the ExchangeEnergyContract class, the list of verification rule objects is created in the form of lambda expressions.It is worth noting that verification rules 4 and 6 are dedicated to processing objects of the TransactionCross and TransactionGrid classes, respectively.The constructor also sets configurations as distinct lists of references to verification rule objects for each transaction type.The checkSC() method, implemented in the smart contract abstract class, operates on the first configuration of verification rules.In a specific smart contract class, this method should be overloaded as many times as there are additional transaction types.In the example considered, two more methods had to be written.One method for the cross-community transaction type and one for the to-grid transaction type.Importantly, if the checkSC() method is called with a transaction type other than declared for the smart contract, it will execute correctly and return a false logical value.The true logical value, which proves the correct verification of the transaction, may be returned only for one of the considered transaction types.
The source code of the checkSC() method for handling TransactionGrid transactions was shown in Listing 4. On invocation of the checkSC() method, Java verifies the parameter type and calls that appropriate method from those overloaded.
The calling of the overloaded checkSC() method is shown in the UML Sequence diagram (Figure 6).The order in which the verification rules appear in the configuration matters, because the checkSC() method evaluates the verification rules in the order they were set in the configuration.Evaluation of the configuration is aborted if a single verification rule is not met.Such a way of evaluation shortens the smart contract checking time.
The TransactionIn and TransactionGrid classes were implemented in the same way as the class for cross-community transaction type.Both classes inherit from the Transaction class and terminate inheritance by employing the final keyword.The use of the Transaction class, in rulesList variable, also allows the processing of a list of rules on each of the specific types of transactions by the same method.
The source code of the AdapT v2.0 design pattern is available in the publicly accessible GitHub repository [43].
Reuse and Redundancy Analysis
The effectiveness of the use of the software source code affects its maintenance.Raising the level of its reuse facilitates modifications and reduces the scope of testing.The author has introduced a measure U sc as the percentage of reused verification rules under a smart contract.The U sc is expressed by Equation (1). where: C-the number of configurations in the smart contract; u r i -the number of reused verification rules in i-th configuration in the smart contract; u i -the number of verification rules in i-th configuration in the smart contract.
The value of the source code efficiency reuse U sc for the smart contract considered in the article was calculated: × 100 = 58.3%.
Configurations were taken for calculations in order from the most numerous to the least numerous.A score above 50% indicates a high level of source code reuse of validation rules.It should be remembered that if the design pattern was not applied, such would be the level of redundancy of verification rules.
At run-time, it is also worth determining the level of efficiency of using the set of verification rule objects in configurations.The author has proposed the measure D sc as the percentage of redundant verification rule objects for a smart contract.The measure can be expressed by Equation (2). where: o c i -number of newly created objects in i-th configuration of the smart contract; v sc -number of unique verification rules in the smart contract.
The value of the percentage of redundant verification rule objects D sc for the smart contract considered in the article was calculated: × 100 = 0%.This means that verification rule objects are fully used by configurations.In addition, making changes in checking between transaction types does not create new objects or drop existing ones.As a result, the garbage collector is not involved in the operation of the software.It should also be underlined that no verification rule object from the verification rule object list is left unused.
In a recently published research paper, Khan et al. [35] showed that the overall cloning rate of Solidity smart contracts is 30.13%, of which 27.03% are exact duplicates.The AdapT pattern employed to design the smart contract allows for achieving a verification rule cloning rate of 0%.Therefore, developing design patterns that increase the reuse level of the source code of smart contracts seems to be one of the appropriate directions of research work.
In addition, the efficiency of the run-time use of operational memory is of great importance.Working with various transaction types may involve the constant instantiation of many extra objects.Especially when transactions of various types are evaluated alternately.The author has proposed the measure of object creation efficiency O sc T as the mean value of the number of verification rule objects created at runtime in the smart contract for a single transaction.The measure can be expressed by Equation (3). where: T-the number of checked transactions; o t i -the number of newly created objects for i-th transaction.Calculated values of O sc T for the considered smart contract for selected quantities of checked transactions were presented in Table 1.The value of the measure of object creation efficiency O sc T decreases for the considered smart contract as the number of checked transactions increases.It stems from the fact that the complete set of verification rule objects is created when a smart contract object is instantiated.All transactions are served by the same collection of objects regardless of the number of transactions.For the smart contract considered in the article, six objects are created, one for each of the verification rules.That set of six verification rule objects verifies any number of transactions.No other objects are created for verification rules.
Performance Tests Results
The purpose of the performance tests was to check how quickly transactions are checked by a smart contract designed following the AdapT pattern.As far as the execution environment is concerned, the tests were carried out on a MacBook Air with the Apple M2 processor, 16 GB Random Access Memory (RAM), and 256 GB Solid State Drive (SSD).The MacBook Air worked under the macOS Sonoma 14.3 operating system.To conduct performance tests, a separate TestContract smart contract testing class was designed.The class contains two methods: conductTest() and runTests().The first method measures the evaluation time of a smart contract.The method uses the System.nanoTime() to obtain an accuracy of one nanosecond of time measurement.The second method is responsible for performing the appropriate number of measurement repetitions.The source code of the TestContract class is shown in Listing 7.
The evaluation time of a specific number of transactions by the smart contract E sc n was adopted as the basic performance measure.Each test has been repeated 50 times to obtain the mean value of E sc n .Tests were conducted for the smart contracts with 3, 4, and 5 verification rules.A total of 450 tests were run.
Figure 7 depicts test results for the number of transactions in the following range, n ∈ ⟨100, 000; 10, 000, 000⟩.The results in the figure are presented on a logarithmic scale.
One of the facts that stem from the results presented in Figure 7 deserves to be emphasized.The evaluation time of 10,000,000 transactions is below 0.25 s, regardless of the considered number of verification rules in the smart contract.Currently, the Solana blockchain framework is regarded as one of the quickest offering a throughput of up to several dozen thousand transactions per second [44].The results obtained illustrate the performance potential of smart contacts designed according to the pattern.
It is also worth visualizing the mean evaluation time of a single transaction.Figure 8 shows the values for this measure calculated from the results obtained when evaluating the considered transaction volumes.Results are shown on a linear scale and are expressed in nanoseconds.The average time to check a single transaction is practically constant and is approximately 25 nanoseconds.This means that the smart contract validation mechanism in the pattern has been designed correctly.The increase in evaluation time is directly proportional to the number of transactions evaluated.
Discussion and Limitations
In the construction of the pattern, two layers have been distinguished: abstract and specific to a particular smart contract.The use of an abstract layer is a recommended software design practice.Recently, Spray et al. [45] have shown the positive impact of the Abstraction Layered Architecture on software reusability and testability.The abstract layer of the AdapT pattern is free of implementation-specific classes or interfaces.The layer consists of two abstract classes and uses only interfaces from Java Standard Edition packages, i.e., List and Predicate.The VerificationRule interface is not needed anymore.Instead, verification rules are handled by the Predicate functional interface.Additionally, the abstract smart contract class is independent of any specific transaction class.Generally, both abstract classes in the abstract layer are in this version of the pattern completely independent of the implementation of a specific smart contract.Thanks to this, those abstract classes are reusable without the need for any changes.
The construction of the concrete layer has also been simplified.Currently, no class is needed for any verification rule in the concrete layer.Verification rules are stored as lambda expressions using the Predicate interface.It is worth adding that they are stored in one variable, which makes maintaining verification rules within the smart contract easier.Such a manner of verification rule design makes smart contracts less prone to errors.The construction of the concrete layer reduces the number of classes and raises software maintainability.
The pattern also employs polymorphism as one of the basic object-oriented paradigms.The transaction verification method the checkSC() was overloaded for the considered transaction types.The usage of overloaded methods simplified the source code and reduced the number of operations needed.The overloaded methods eliminated conditional statements from the implementation of the verification mechanism invocation.
Java was used because it is a general-purpose language.The development community of this language is large.Additionally, there are available open-source components written in Java.The language itself allows the use of both object-oriented and functional programming structures.Java is also used to program smart contracts on Corda, Hyperledger Fabric, IBM Blockchain, Ethereum, and Neo.
The pattern reduces the redundancy of verification rule objects at run-time to zero.Verification rules are reused among configurations within the single smart contract.The possibility of reusing verification rules between smart contracts was not considered.The paper deals with configurations.In the previous version of the pattern, there were only two possible configurations for two transaction types.Now the design allows for flexibly adding and handling any number of configurations.Configurations may find many applications.Configurations may be applied to handle on-chain and off-chain transactions realized by the same smart contract.Blockchain is also used for authentication and configurations of a smart contract can be used to authenticate various roles.Additionally, the idea of the reuse of verification rules may be adopted for logical conditions in various methods of a smart contract.This should generally increase the maintainability of smart contract source code.
The accuracy of the obtained smart contract execution time is at the level of 1 nanosecond.Because the evaluation time of a single transaction is approximately twenty-five nanoseconds, the measurement precision is sufficient.Especially since the aim was to show the behavior of the smart contract evaluation mechanism with significant transaction volumes.
Conclusions
The article presents a smart contract design pattern that allows for handling many types of transactions.The paper contains both the design of the pattern and its implementation in Java language.The pattern's structure combines object-oriented and functional programming mechanisms.The adoption of sealed classes increases the security of processed transactions.The Predicate<T> functional interface and lambda expressions were employed to reduce the number of classes.As a result, software maintenance costs are also lowered.The design of the pattern allows the verification rules to be reused within the smart contract.The usage of the pattern ensures that the redundancy ratio of verification rule objects at runtime is zero.It also allows for the effective use of operating memory by the smart contract.The implementation of the pattern is independent of the blockchain platform.As a result, it was possible to conduct performance tests of the written software independent of other elements of the blockchain technology, in particular the consensus algorithm.The tests were carried out for a wide range of the number of processed transactions.The evaluation time of a smart contract for <100,000; 10,000,000> transactions takes between <0.0025; 0.25> seconds.Importantly, the analysis of the tests showed that the increase in the evaluation time is directly proportional to the number of transactions checked.The performance test results clearly show that smart contracts as software have great potential for processing large volumes of transactions, far beyond the capabilities of currently available environments.
In the context of further work, Rust language is gaining increasing attention in the community.A general-purpose language may have various applications.Blockchain is one of them.Rust enforces memory safety and a restrictive model of data ownership.The language eliminates the need for a garbage collector.Additionally, Rust supports
Figure 1 .
Figure 1.Configurations share verification rule objects from the same list.
Figure 2 presents a UML Class diagram with both classes in the Abstract layer of the AdapT pattern.
Figure 2 .
Figure 2. Classes in the Abstract layer of the AdapT v2.0 pattern.
Listing 2 .
p u b l i c a b s t r a c t c l a s s A b s t r a c t S m a r t C o n t r a c t { p r o t e c t e d L i s t < P r e d i c a t e < A b s t r a c t T r a n s a c t i o n >> r u l e s L i s t = new A r r a y L i s t < >() ; p r o t e c t e d L i s t < L i s t < P r e d i c a t e < A b s t r a c t T r a n s a c t i o n >>> c o n f i g u r a t i o n s = new A r r a y L i s t < >() ; p u b l i c A b s t r a c t S m a r t C o n t r a c t ( ) { c o n f i g u r a t i o n s .add ( new A r r a y L i s t < >() ) ; } p u b l i c boolean checkSC ( A b s t r a c t T r a n s a c t i o n t r ) { boolean c o r r e c t = f a l s e ; f o r ( P r e d i c a t e < A b s t r a c t T r a n s a c t i o n > vR : c o n f i g u r a t i o n s .g e t ( 0 ) ) { c o r r e c t = vR .t e s t ( t r ) ; i f ( ! c o r r e c t ) break ; } r e t u r n c o r r e c t ; } } The source code of the AbstractTransaction class is shown in Listing 2. The source code of the AbstractTransaction class.package adapt ; p u b l i c a b s t r a c t c l a s s A b s t r a c t T r a n s a c t i o n { }
Figure 3 .
Figure 3.The SendEnergy use case with various external users.The action performed in the use case is the same, but the conditions checked differ depending on the transaction type.The example assumes the following set of verification rules used by the three transaction types considered: • The source of the transaction must be different from the target of the transaction; • Energy quantity to transfer must be greater than zero; • Energy surplus in the source node must be greater or equal energy quantity to transfer; • Source community must differ from target community; • Target need must be greater or equal to energy quantity to transfer; • The target is the subnet energy grid.Using the AdapT pattern, the Send energy use case can be implemented with one smart contract.Figure 4 depicts a UML Class diagram with classes that both constitute the abstract layer of the AdapT pattern and implement classes of the concrete ExchangeEnergy smart contract.
Figure 4 .
Figure 4.The AdapT pattern with abstract and concrete classes.
Figure 5 .
Figure 5. Classes for transaction types for the SendEnergy use case.
Listing 4 .
The source code of the checkSC() method for to-grid transactions.p u b l i c boolean checkSC ( T r a n s a c t i o n G r i d t r ) { boolean c o r r e c t = f a l s e ; f o r ( P r e d i c a t e < A b s t r a c t T r a n s a c t i o n > vR : c o n f i g u r a t i o n s .g e t ( 1 ) ) { c o r r e c t = vR .t e s t ( t r ) ; i f ( ! c o r r e c t ) break ; } r e t u r n c o r r e c t ; }
Figure 6 .
Figure 6.Calling the smart contract verification method.
Table 1 .
Values of O scT for the considered smart contract. | 9,255.4 | 2024-03-07T00:00:00.000 | [
"Computer Science"
] |
Fredholm pseudo-gradients for the action functional on a sub-manifold of dual Legendrian curves of a three dimensional contact manifold ( M 3 , α)
We prove in this paper that the intersection numbers between periodic orbits have an intrinsic meaning for the variational problem ( J , C β ) {Bahri (Pseudo-Orbits of Contact Forms Pitman Research Notes in Mathematics Series No. 173, 1984), Bahri (C R Acad Sci Paris 299, Serie I 15:757–760, 1984), Bahri (Classical and Quantic periodic motions of multiply polarized spin-manifolds. Pitman Research Notes in Mathematics Series No. 378, 1998)}, corresponding to the periodic orbit problem on a sub-manifold of the loop space of a three dimensional compact contact manifold ( M , α) .
To a certain extent, these two difficulties are intertwined. In many problems of Conformal Geometry, e.g. the Yamabe and related problems, the associated variational problems [1] are (locally) Fredholm, but they do not verify the Palais-Smale condition. Suitable techniques [9][10][11]27] have then been developed to overcome, at least partially, this difficulty.
In the area of Hamiltonian Systems, the Fredholm assumption and the so-called (P.S) condition are "easy" (in that they are now classical) to verify for Lagrangian formulations (e.g. [12,17,28] including brake-orbits [33]). In the new formulation developed by P.H. Rabinowitz [25] in 1978, with the introduction of the action functional 1 0 p iqi on the space H 1 2 (S 1 , R 2n ), both conditions are verified for example through a Galerkin approximation by finite dimensional spaces.
This framework has also been used by C. Conley and E. Zehnder [13] for the solution of the Arnold conjecture on tori. However, as mathematicians moved away from the R 2n -framework and tried to solve the Arnold conjecture in full generality or tried to solve the Weinstein conjecture [32] 1 , they found themselves without an appropriate variational formulation for the periodic orbits problem for contact vector-fields.
It is not an easy task, even in the framework of cotangent bundles of finite dimensional manifolds, for Hamiltonians that are not convex in the momentum variables (no Lagrangian formulation), to define a Fredholm framework for this problem, see e.g. [21]. The space H 1 2 (S 1 , M 2n ), which is the natural space (e.g. in a symplectic formulation) for the action functional, is not well-defined because H 1 2 (S 1 , R 2n ) does not embed in L ∞ .
Several methods have been devised to overcome this difficulty. For example, A. Floer [16], using the pseudo-holomorphic framework introduced by M. Gromov [19], was successful in extending the results of C.
Conley and E. Zehnder [13] to the framework of compact symplectic manifolds. Also, in the contact framework, H. Hofer [20] was able, using this pseudo-holomorphic framework and the construction of a special disk for over-twisted contact structures, to prove the existence of one periodic orbit for the related contact vectorfield (and did in this way solve positively, to a great extent, the three-dimensional version of the Weinstein conjecture).
However, despite this progress, the full understanding of the Morse relations between the periodic orbits of a given contact vector-field could not be achieved.
For example, with pseudo-holomorphic curves (assuming their existence, a non-trivial matter), one can try to understand these Morse relations through the moduli spaces of such curves [15]. However, along deformations of contact forms, these moduli spaces are not stable. There are "blow-ups", with discontinuities in the Fredholm index and failure of compactness. Beyond the issue of existence of moduli spaces of pseudo-holomorphic curves, one finds himself facing again the two fundamental difficulties described above.
Very early, we have defined, in collaboration with D. Bennequin [3], a variational framework for the periodic orbits problem for contact vector-fields on a three-dimensional closed and compact contact manifold (M 3 , α). In this variational framework see e.g. [2,5,6], the action functional J (x) = 1 0 α(ẋ) was studied on a sub-manifold C β = {x ∈ H 1 (S 1 , M); β(ẋ) = dα(v,ẋ) = 0; α(ẋ) = a} of the loop space of M 3 . a, in the definition of C β , is a positive constant that is not prescribed, it varies with the curve x(t); v is a non-singular vector-field in kerα, β verifies the condition (A) : β is a contact form with the same orientation than α, see Sect. 6 below for a considerable weakening of this condition.
Very early on also [2], we had noted that this variational problem failed both the Fredholm assumption and the Palais-Smale condition.
We have overcome, in various (different) ways the second difficulty in our work, see [5][6][7][8] in particular. However, we could never overcome the violation of the Fredholm assumption, although we did reduce it in [8] to a violation of this assumption at the periodic orbits themselves. We were able in [8] to formulate a simple condition. Under this condition and for a special pseudo-gradient, see [8], the intersection operators ∂ per and ∂ ∞ do not mix in between creations and cancellations of periodic orbits.
We prove in the present paper that there is a pseudo-gradient flow for (J, C β ), that can be continuously tracked along deformations of contact forms, for which the Fredholm assumption at the periodic orbits is verified (after [8], this is all what is needed) and that, for this pseudo-gradient, the intersection number between two periodic orbits of consecutive indexes is defined intrinsically (as described above, in the compact, finite dimensional framework). Accordingly, with the use of this flow and the additional work in [5,6] and [8], the variational problem (J, C β ) becomes a "Fredholm framework" for the finding of periodic orbits to the contact vector-field ξ of α (β = dα (v, .), v ∈ kerα, see [2,4]), "stable" under deformation. This is already a significant progress in the effort to find an appropriate framework for the problem of periodic orbits. However, further progress is much needed to extend these techniques to higher dimension and to entirely remove conditions (A) and (A) t .
We do not claim here to have the final framework for this kind of variational problems. The present paper rather asserts a direction of research and states positive (global) results of existence related to this direction. This is a short paper and its main result is stated in Proposition 2.1. This Proposition is about the intersection number of two periodic orbits when there are no other periodic orbit or critical point at infinity in between their energy levels (for the action functional). This can be readily extended to allow for intermediate critical points with zero intersection numbers with the dominating or with the dominated periodic orbit, depending on their index, see above, at the beginning of this Introduction.
The proof of this Proposition 2.1 assumes the knowledge of the results of [5,6,8]. We proceed now with our precise claims and proofs: Let α t be a deformation of contact forms on a contact closed manifold M 3 and let v t be a family of continuously varying vector-fields in their kernel (α t (v t ) = 0). Let us assume that the condition ) is a contact form with the same orientation than α t is verified all along the deformation. We will indicate at the end of this paper how to get rid of this condition.
As in [2,4,5], t 2m , which we also denote 2m , is the space of curves made of mξ t -pieces of orbits alternating with m ± v t -pieces of orbits. ξ t is the Reeb vector-field of α t .
Let a b be two values such that J t has no critical point at infinity in (J t ) −1 ([a, b]) but for the (δ (m) +w) ∞ maybe (these (δ (m) + w) ∞ are the critical points at infinity built with "Dirac masses", i.e. back or forth or forth and back runs along v, above some point of the periodic orbit w, see [5] and [8] for more details, [8], Appendix 1 in particular). Assume furthermore that the deformation α t has been "adjusted", using the techniques of [5], p 85-93, e.g. Proposition 15, see also [6], p 473-474 for an earlier use for this proposition to "adjust" the v-rotation along a simple periodic orbit, so that: (i) Every w t 2m+1 such that a J t (w t 2m+1 ) b is a simple elliptic periodic orbit; whereas every w t 2 p such that a J t (w t 2 p ) b is a simple hyperbolic periodic orbit. w t 2m+1 has Morse index (2m + 1), w t 2 p has Morse index 2 p. (ii) Given two periodic orbits w t 2k+1 and w t 2k in (J t ) −1 ([a, b]), of Morse index (2k + 1) and 2k respectively, we assume that either a cancellation (w t 2k+1 /w t 2k ) occurs at the time t = t 0 ; or that the level J t (w t 2k+1 ) crosses the level J t (w t 2k ) at the time t = t 0 . α t is then "adjusted", if needed, so that the v-rotation along the simple elliptic periodic orbit w t 0 2k+1 is 2kπ + θ, θ ∈ (0, π). (iii) On the other hand, if instead of a cancellation/crossing (w t 2k+1 /w t 2k ), a cancellation/crossing , at the time t = t 0 , the v-rotation along the simple elliptic periodic orbit w t 0 2k−1 is (2k − 1)π + θ, θ ∈ (0, π). In addition, we assume that the v-rotation along w t 2k , starting from any point along w t 2k−1 , is 2kπ + o(π), just as in Section 4 of [8]. There is no loss of generality in assuming that (i)-(ii)-(iii) holds, see Proposition 15 of [5] and [6], p 473-474. In the case of cancellations, these conditions are verified for a deformation α t in general position, without the need for any further adjustment.
After (i)-(ii)-(iii), we claim that the following holds in that is "Fredholm" or "symplectic" ([8], Definition 1): no tangency between W s (w t r ) and W t s (δ + w t r −1 ) ∞ ) occurs over the deformation for any two periodic orbits w t r , w t r −1 , of respective indexes r, . The family Z t varies in a differentiable way with t and defines a flow on each 2m , m ∈ N .
Proof of Proposition 2.1 In all our arguments below, the "Dirac masses" built over the various curves in the deformation process (on the stable and unstable manifolds of the various critical points and critical points at infinity involved in the arguments) are suitably approximated with back and forth or forth and back runs along the vector-field v, these runs along v being separated by tiny ξ -pieces that eventually become larger as the two ±v-jumps become small and are "pushed away" one from the other one. It is to this set of approximating curves, rather than to the "infinitely" contracted curves with "Dirac masses", that the arguments for elliptic orbits w 2k−1 and hyperbolic periodic orbits w 2k are applied below: the infinitely contracted "Dirac masses" could otherwise resolve themselves into two confounded zero ±v-jumps through the "pushing away" and "widening process" of [6] and the arguments for Proposition 2.1 would then become less transparent.
We start now the proof of Proposition 2.1: since there are no critical point at infinity in is modelized, see Proposition 2.1, p 469 of [6], with m single ±v-jumps that can be tracked over decreasing flow-lines.
with curves that support 2k simple ±v-jumps separated by ξ -pieces of orbits.
If one of these flow-lines enters an L ∞ -neighborhood (in graph) that is small enough of (δ + w t 2k−1 ) ∞ , then the curves on this portion of flow-line must have at least two non-zero ±v-jumps. Using the arguments of Section 3 of [8], "Bypassing a simple elliptic periodic orbit", such a flow-line will never end at w t 2k−1 .
are thereby forbidden with such a flow and the intersection number i(w t 2k , w t 2k−1 ) does not change with t. In addition, all these flows can be deformed one into the other, those that do not introduce any companions to existing single ±v-jumps in (J t ) −1 ([a, b]) as well. For all of these, the flow-lines that enter an L ∞neighborhood (in graph) that is small enough of (δ + w t 2k−1 ) ∞ do not abut at w t 2k−1 later. As we deform continuously pseudo-gradients over the time of the deformation, we find a flow in (J t ) −1 ([a, b]) which we may assume to not introduce companions to existing simple ±v-jumps in this "energy slice".
If we continuously deform this flow, among pseudo-gradients that have the same property in All pseudo-gradients with this property can be deformed one into the other. Among these, there is a "compact" pseudo-gradient which is almost explicit on W u (w t 2k ) as the two energy levels, the one of w t 2k and the one of w t 2k−1 become closer and closer. The intersection number can be computed on this compact one. The claim of the Proposition 2.1 follows in this case (see some further precisions below, when considering configurations such that the ±v-jumps of the "Dirac mass" are not well-defined).
A similar phenomenon occurs for a pair w t ), but the proof is different. Again, the curves on the flow-lines out of w t 2k+1 support (2k +1) simple ±v-jumps that can be continuously tracked. If a flow-line enters a small L ∞ -neighborhood of a (δ + w t 2k ) ∞ , then two of its ±v-jumps are large. We also observe that we can take them to be consecutive ±v-jumps.
When this flow-line reaches a small L ∞ -neighborhood of w t 2k , the behavior of the related configurations can be understood as follows: the two consecutive ±v-jumps are still non-zero, but small. Completing the "widening process", see [6], Proposition 20, p 518, between these two ±v-jumps, we can bring the v-rotation on the ξ -piece separating them to be π − π 2k+1 + o(e −k ) as in Section 1 of [8]. Their orientations have not been reversed and they are still non-zero ±v-jumps at the end of this process. They are now "locked in their positions". The remaining (2k − 1) ± v-jumps have to be "rearranged" through the process of "pushing away" so that the v-rotation on any ξ -piece separating two consecutive ±v-jumps is i(w t 2k , w t 2k−1 ). The v-rotation on the "external" nearly ξ -piece (it is "broken" with the remaining (2k − 1) ± v-jumps. All ±v-jumps are assumed, without loss of generality to be 0(e −k 2 )) separating the two non-zero consecutive ±v-jumps that were involved in the formation of the "Dirac mass" is 2k(π − π 2k+1 + o(e −k )). It follows in particular that this second step in the re-arrangement process can be completed so that the two consecutive ±v-jumps involved in the formation of the "Dirac mass" do not change location.
We therefore need to understand better now the first part of this process as the "Dirac masses" build up and are still large. We have developed such an understanding in [8], Appendix 1. However, in our present line of arguments, we seek an understanding of these phenomena when no companions are introduced to existing ±v-jumps, although we will be also indicating the modifications needed in the case companions are introduced as in [6] and [7]. Consequently, we also need to adjust to this framework our understanding of the process of formation of "Dirac masses".
Zoology of "Dirac masses"
A slightly different "zoology" for these "Dirac masses" holds in this new context. We now describe this zoology: "Dirac masses" contain a back and forth or forth and back run along v. Let us assume, without loss of generality, that we are in the latter case.
The function: x s is the v-orbit through x 0 , is relevant to the formation of the "Dirac mass", see [5], pp 28-29 and [8], Appendix 1 for more details.
A "Dirac mass" of the type indicated above can be built whenever θ(s) is negative for some positive s, see [8], Appendix 1 also.
However, depending on the behavior of θ(s) for s ∈ (0, ∞), we may encounter different configurations involving different outcomes. θ(s), on sub-intervals of [0, ∞), can behave in four basic ways, and also in a fifth way. These are best described with the drawings below: I, II and III can be modified so that they will include more bumps, V also. However, more relevant to V is the behavior of θ(s) when the base point on x 2k varies. In general, V will break into: However, in the case of circle-bundles along v, if xs is x 0 , then this behavior (the one described in V ) survives the change.
Index at infinity of "Dirac masses": In all the cases that we are considering, w m with the addition of a "Dirac mass" may be viewed as a curve of 4 , with one ξ -piece reduced to zero, that is it can be viewed as a critical point at infinity of index i ∞ (at infinity, in 4 ) equal to 0 or to 1. Along deformations, there is an additional parameter and the index at infinity can be equal to 2. Indeed the flow-lines of 4 out of an x m+1 build a stratified space of top dimension 2. These flow-lines must dominate the critical point at infinity defined by this "Dirac mass" and this implies the conclusion.
More specifics about the "zoology", "energy levels": The precise value that the function θ(s) takes at the edge of the "Dirac mass" is irrelevant. All these curves are at the same energy level for J ∞ . This energy level depends only on the base curve w 2k or x upon which the "Dirac mass" is built.
However, the fact that θ(s) is positive or negative at the edges of the "Dirac mass" matters: if θ(s) is negative at the upper-edge, then, see Appendix 1 of [8], a small ξ -piece can be inserted at the "top" of this "Dirac mass" and J ∞ decreases substantially along this process.
Assuming now that θ(s) is non-negative at this upper-edge, we can make the "Dirac mass" longer or shorter. The first zeros for the function θ(s) that we encounter in either direction matter then: as soon as the edge of the "Dirac mass" enters an interval where θ is negative, the process of insertion described above can be completed and J ∞ decreases substantially.
This allows to understand better the behavior of the unstable manifold at infinity (recall that i ∞ = 0 or i ∞ = 1, see above and [8], if the "Dirac mass" is dominated by flow-lines of W u (w 2k+1 ) ∩ 4 ).
If we are to discriminate between the "energy levels" defined by J ∞ for the curves of 4 built with "Dirac masses" as above, then we can define a flow that follows the behavior of the function θ . Then, a "Dirac mass" such as I or III is higher than II and a "Dirac mass" such as V (a) is higher than V (b).
Once this flow on the "Dirac masses" is defined in 4 , the critical points at infinity of index i ∞ = 0 or i ∞ = 1 -they all turn out to be of index i ∞ = 1-are either isolated curves of the type V (b) (their precise value depend on the full definition of the flow); or starting from "Dirac masses" of type I and following the analysis of Appendix 1 of [8], the flow-lines end at "Dirac masses" located at precise points x i 0 such that the upper-edge of the "Dirac mass" verifies: Both types of curves are of index i ∞ = 1. For II, the exit sets may be read on the drawing: The "Dirac mass" can be made longer and shorter so that the function θ is negative at its edge and, see Appendix 1 of [8], a small ξ -piece can be inserted at its top in a J ∞ -decreasing process.
These curves do not therefore dominate w 2k at infinity. The arguments of Section 3 of [8] may therefore be applied to them. They fit in the framework of the critical points at infinity of [8].
If, instead of the arguments of [8], we allow the introduction of companions, as in [5] and [6], we can always "spare" the negative of the positive v-jump of the "Dirac mass" and choose not to introduce companions to one of them.
Introducing companions to the other one, the "Dirac masses" of type II do not appear since it is possible then to insert, there where θ is negative along the positive or the negative v-jump of the "Dirac mass", a small ξ -piece and J ∞ will thereby decrease substantially. The arguments for the expansion of J ∞ along this insertion can be found in [5] and [6].
Let us furthermore observe, along the same line of arguments that, whether the introduction of companions is allowed or not, critical points at infinity of type (III) do not appear: decreasing the size of the "Dirac mass", we reach an interval where θ(s) is negative. Inserting then a small ξ -piece, J ∞ decreases substantially.
For curves of type IV, we can make the "Dirac mass" shorter until it disappears: there is no exit set related to the behavior of θ(s) for s ∈ (0, ∞).
Thus, through various arguments, depending on whether we are using the techniques of [7] or we are using the techniques of [8], critical points at infinity of type (III) either do not appear or can be treated as in Section 4 of [8].
Curves of type I and end of proof of Proposition 2.1
We are left with curves of type I.
Using the arguments of [5] and [8], they are viewed as critical points at infinity. Using companions, see [7], Section 8, they can be bypassed.
If we do not allow the use of companions, we still claim that no flow-line of W u (w 2k+1 , coming out of a neighborhood of a "Dirac mass" of type I, located at an x 0 along w 2k satisfying ( * * ), will abut at w 2k .
To see this, we resume our rearrangement/reordering argument above. The positive v-jump of the "Dirac mass" is locked at x 0 along this process (it has been chosen as γ , see Section 3 of [8], in our argument). Indeed, this positive v-jump does not change location as the "Dirac mass" decreases in size along the unstable manifold. This holds also, as has been pointed out above, when the curves support several other small ±v-jumps that are used to modelize the H 1 0 -unstable manifold of this critical point at infinity.
We now re-scale the v-rotation along w 2k so that, all along the deformation, these points x 0 (there might be several such x 0 s) are located in two consecutive intervals of "positive type": these are two consecutive intervals along which J " ∞ (w 2k ) is positive along a curve defined with the insertion of a small ±v-jump located at some point along x 2k , see Section 1 of [8]. Such "positive" intervals are separated by intervals where J " ∞ (w 2k ) is negative along the same type of directions.
These two consecutive intervals of positivity for J " ∞ (w 2k ) are evolved over the deformation of contact forms so that, at any time t along this deformation, all the x 0 s (with i ∞ = 1) such that ( * * ) holds for some s ∈ R − {0} are included in one of these intervals. This involves of course a continuous rescaling of the v-rotation along w 2k that can be completed as in [5], pp 85-93, Proposition 15 in particular.
In this way, we can include in intervals of positivity for J " ∞ (w 2k ) all the base points of the "Dirac masses" such that their top level is not a local maximum in the set of "Dirac masses". Indeed, any "Dirac mass", as explained above, must be of index at least 1, this is embedded in the construction and a direction of negativity can be recognized as living along the "Dirac mass", as it changes size.
For these "Dirac masses", as we re-arrange the configuration when the ±v-jumps become small, the positive v-jump of the "Dirac mass" can be kept locked at x 0 . For configurations coming from "Dirac masses", the rearrangement can be completed by pushing all the ±v-jumps away from the two ±v-jumps of the "Dirac mass, then "pushing the negative ±v-jump away from the positive one and adjusting then the rotation. Whatever happens does not change one fact: one ±v-jump besides the positive v-jump of the "Dirac mass" remains non-zero, so that, after Section 3 of [8], if we choose the positive direction E + at this configuration to be modeled by a ±v-jump located at x 0 (it can be done whenever x 0 is in a positivity interval for J " ∞ (w 2k )), this configuration is not in the stable direction for w 2k and the configuration can be moved down, past w 2k .
We are left with the "Dirac masses" of higher index at infinity that have their base point located in an interval where J " ∞ (w 2k ) is non-positive. If the base point is in an interval on negativity, then the index at infinity of the "Dirac mass" is at least 3: two directions of negativity are one along the "Dirac mass" as explained above, the other one because this "Dirac mass" top level is a local maximum among the top levels of the "Dirac masses". The third direction comes the possibility of changing the relative sizes of the positive and the negative large ±v-jumps of the "Dirac mass", thereby creating a third negative direction since the base point is an interval of negativity for J " ∞ (w 2k ).
Such curves of 4 cannot be dominated along the deformation, the index at infinity is too large. Another ±v-jump must be non-zero. The argument of Section 3 of [8], with the three edge rule can be applied and work without the need for Section 11 of [8] because γ 0 may be chosen as one of the two (initially) large ±v-jumps of the "Dirac mass" since these two follow each other with opposite orientations: since these two ±v-jumps are so close, rearrangement can be completed by "pushing away" all the other ±v-jumps from this pair and never "pushing" one of them away from these (2k − 1) other ±v-jumps. In addition, the final rearrangement, with these two special ±v-jumps finding their final position can be completed with the use of the "widening process", see [6], between them, so that their respective orientations is never altered in this process. It can only be altered by the fact that we "pushed away" the other (2k − 1) ± v-jumps from them. Repetitions are then preserved and the rearrangements around either choice are the same. A choice for the positive direction at a configuration among these can be completed in a compatible way over the switch of choices for γ 0 amongst these two ±v-jumps, see Section 3 and Section 11 of [8]. It follows that these configurations are also moved down, past w 2k (some further precisions are needed and provided below, when the two large ±v-jumps of the "Dirac mass" are not welldefined).
At specific times, "Dirac masses" cancel themselves topologically; that is a "Dirac mass" whose top level is critical among top levels, but is not a local maximum cancels with one whose top level is a maximum. Since we are requiring that the first species have all their base point in intervals of type E + , we have to allow for the second species to cross over, at certain times, from E − to E + , before cancellation.
At these specific times along the deformation, such a "Dirac mass" is located at a node, moving from E − into E + or vice-versa. It has at least two non-zero decreasing directions at infinity in 4 , one as all "Dirac masses" do have; it is related to their length; the other one is related to its top level.
Flow-lines out of w 2k+1 in 4 are of dimension 2. Adding the deformation parameter, we find a set of dimension 3. Such a set cannot dominate a critical point at infinity-such as the above "Dirac mass"-of index at infinity larger than or equal to 2 but at specific times along the deformation. We need to warrant that these specific times are not the times at which these "Dirac masses" cross E 0 . This amounts to check that, as these "Dirac masses" cross E 0 , we can still perturb, far away, near w 2k+1 , the deformation and our flow-lines in 4 so that they do not dominate these "Dirac masses".
The argument is standard as the value of the contact form near w 2k+1 is very much independent from the v-rotation on the hyperbolic orbit w 2k . The claim follows.
More complicated configurations and the verification of the Fredholm assumption
The arguments provided above rule out the violation of the Fredholm assumption for configurations where the large ±v-jumps of the "Dirac mass" are well-defined. We prove that the arguments extend to the general case: Assuming that we are considering here a Morse relation as above between a w 2k and an elliptic orbit w 2k−1 and assume that the positive or the negative edge of the "Dirac mass" over this periodic orbit is not well-defined, i.e. two or more than two * s define it or there is a tiny or zero ±vjump in between the two large edges of the "Dirac mass". We will consider the case of one zero ±vjump in between these two large ±v-jumps. The other cases force the occurrence of more repetitions, at least two as the ξ -piece in between these two large ±v-jumps is tiny and does not support any Under such an occurrence, there is a forced repetition in between the 2k ± v-jumps of the configuration. Two ±v-jumps to the least are non-zero and their orientations force the existence of a repetition in between them.
The use of any * among the 2k available ones as a γ will not change this fact along such configurations (repetition and two non zero ±v-jumps). They can be moved down with any such γ and this deformation convex-combines with the decreasing deformation centered at the positive or at the negative (now well defined) ±v-jump of the "Dirac mass" used over the remainder of the set of configurations.
For a hyperbolic orbit w 2k , the configurations out of w 2k+1 having one zero ±v-jump in between the two large ±v-jumps of the "Dirac mass" correspond to a stratified set T of top dimension 2k (in 4k ).
Along a deformation of contact forms, this stratified set might undergo tangencies with the stable manifold of a "Dirac mass" D, assuming that the dimension of this stable manifold is 2k.
Outside of the two large ±v-jumps of the "Dirac mass", (2k − 2) ± v-jumps are available and they can be used, see Section 11 of [8], to build either the stable or the unstable H 1 0 -manifold of the "Dirac mass". They provide at most (2k − 2) unstable directions.
The missing two dimensions in the co-index are related to the curve formed in 4 by the two large ±vjumps of the "Dirac mass". This co-index must then be 2 for D and therefore the index must be two as well. It follows that the index of D is 2k and W u (D) is achieved in 4k (the additional missing ±v-jump is zero).
Considering this "Dirac mass" in 4k+2 , that is adding the zero or the nearly zero ±v-jump in between the two large ±v-jumps of the "Dirac mass", we easily see that it does not provide any additional index since this ξ -piece is very small, tiny in between two large ±v-jumps.
Therefore the total index of this "Dirac mass" in 4k+2 is also 2k. The hyperbolic orbit itself has an index equal to 2k and has, see Section 11, sub-section on Hyperbolic Periodic Orbits of [8], an unstable manifold of dimension 2k in 4k+2 , 2k in 4k as well.
As tangency takes place between T and W s (D), W u (D) lives in 4k and is of dimension 2k. W s (w 2k ) ∩ 4k is of dimension 2k. Therefore, tangency between W u (D) and W s (w 2k ) occurs (in 4k ) at special times that can be made different from the times at which tangency occurs between T and W s (D). The conclusion follows.
If instead of a tangency as above, we have a domination, then the dominated chain is of dimension (2k − 1) at most, (2k − 2) transversally to the flow and the arguments used in Section 11 of [8] work over transitions, switches in γ s etc., see [8] for more details.
Observe that a "Dirac mass" of index (2k + 1) does not have the "Dirac Mass" D in its boundary. Indeed, the additional ±v-jump that is zero for D cannot be outside the two large ±v-jumps of the dominating "Dirac mass" over the flow-lines of the domination: being small, it can be "pushed away" from them and it will never "enter" between them.
If this ±v-jump is one of two sizable parts of a large ±v-jump of the "Dirac mass" D, then the flow-lines that come to this configuration in 6 can be seen to come from a level much higher than the level of the periodic orbit w 2k and therefore, they do not come from a "Dirac mass" associated to this periodic orbit since the energy level of such a "Dirac mass" is very close to the energy level of the corresponding periodic orbit.
It remains to study the case when this additional ±v-jump is also in between the two large ±v-jumps of the dominating "Dirac mass". Computing the H 1 0 -index of such a "Dirac mass", we find at most (2k − 2), since there are at most (2k − 2) ± v-jumps outside of the two large ±v-jumps of the "Dirac mass". To reach (2k + 1), we would need that the critical point of 4 corresponding to the two large ±v-jumps of this "Dirac mass" to be of index 3 in 4 . This is different from D and therefore the two tangencies W u (w 2k+1 ) with the stable manifold of the dominating "Dirac mass" on one hand and the tangency of T with the stable manifold of D do not occur at the same time by general position arguments.
Therefore, we may assume that this additional ±v-jump is not in between the two large ±v-jumps of this dominating "Dirac mass". To reach D, this ±v-jump would have to travel along a decreasing flow-line and that is ruled out by the previous arguments.
Let us observe, to conclude our argument that we can rule out as follows the case of double-tangency, T with W s (D) on one hand and W u (w 2k+1 ) with the stable manifold of a "Dirac mass" of index (2k + 1) occurring together, over the same process. Indeed, then, the additional ±v-jump is in between the two large ones and therefore it does not provide any additional index so that it is not possible to have a double-tangency.
Outline for the removal of condition ( A) t
The arguments used in [7] and [8], use (A) t in one basic fact: the unstable manifold for a simple periodic orbit w m is achieved in the space 2m .
This holds true under the weaker assumption that (A) t is verified in a neighborhood of w m and this, in turn, holds true-after rescaling the v-rotation along w m using the techniques of [5]-if there is a globally defined v in kerα such that its total rotation around w m in a ξ -transported frame is positive.
If this does not hold and the total rotation for any globally defined, non-singular, v in kerα is negative, then one might have to modify the functional and use − 1 0 α(ẋ)dt near w m or use the same functional, but increase it instead of decreasing it. This would require further work.
However, under the assumption that (A) t is verified near w m , the arguments of [7,8] can be carried out with one additional difficulty: the spaces C β t and the t 2m might have singularities. We have already understood the location of these singularities [4], p 19 for C β , not yet for the 2m s. However, we have not yet built a decreasing deformation for our variational problem through these singularities.
Would this be achieved, the homology would extend under a quite weakened version of (A) t . This condition might be entirely removed after a modification of the functional, starting as indicated above, but this would require more work.
It is worth mentioning here that we can always assume that, at a given periodic orbit, the rotation of υ is monotone, either positive or negative. Accordingly, one finds two Morse complexes; one for the positive rotation, as above and [8]. The other one is related to the functional −J (x) = − 1 0 α x (ẋ) and to the periodic orbits of the second type. The two Morse complexes are, by an argument of general position, independent of each other. Along this line, the present results and the results of [8] can be generalized.
Open Access This article is distributed under the terms of the Creative Commons Attribution License which permits any use, distribution, and reproduction in any medium, provided the original author(s) and the source are credited. | 9,077.6 | 2014-02-15T00:00:00.000 | [
"Mathematics"
] |
Bone-Metabolism-Related Serum microRNAs to Diagnose Osteoporosis in Middle-Aged and Elderly Women
Abstract Objective: Postmenopausal osteoporosis (PMOP), a chronic systemic metabolic disease prevalent in middle-aged and elderly women, heavily relies on bone mineral density (BMD) measurement as the diagnostic indicator. In this study, we investigated serum microRNAs (miRNAs) as a possible screening tool for PMOP. Methods: This investigation recruited 83 eligible participants from 795 community-dwelling postmenopausal women between June 2020 and August 2021. The miRNA expression profiles in the serum of PMOP patients were evaluated via miRNA microarray (six PMOP patients and four postmenopausal women without osteoporosis (n-PMOP) as controls). Subsequently, results were verified in independent sample sets (47 PMOP patients and 26 n-PMOP controls) using quantitative real-time PCR. In addition, the target genes and main functions of the differentially expressed miRNAs were explored by bioinformatics analysis. Results: Four highly expressed miRNAs in the serum of patients (hsa-miR-144-5p, hsa-miR-506-3p, hsa-miR-8068, and hsa-miR-6851-3p) showed acceptable disease-independent discrimination performance (area under the curve range: 0.747–0.902) in the training set and verification set, outperforming traditional bone turnover markers. Among four key miRNAs, hsa-miR-144-5p is the only one that can simultaneously predict changes in BMD in lumbar spine 1–4, total hip, and femoral neck (β = −0.265, p = 0.022; β = −0.301, p = 0.005; and β = −0.324, p = 0.003, respectively). Bioinformatics analysis suggested that the differentially expressed miRNAs were targeted mainly to YY1, VIM, and YWHAE genes, which are extensively involved in bone metabolism processes. Conclusions: Bone-metabolism-related serum miRNAs, such as hsa-miR-144-5p, hsa-miR-506-3p, hsa-miR-8068, and hsa-miR-6851-3p, can be used as novel biomarkers for PMOP diagnosis independent of radiological findings and traditional bone turnover markers. Further study of these miRNAs and their target genes may provide new insights into the epigenetic regulatory mechanisms of the onset and progression of the disease.
Introduction
Postmenopausal osteoporosis (PMOP), which is caused by estrogen withdrawal, is the most frequent type of primary osteoporosis and threatens nearly half of middle-aged and elderly women worldwide [1,2]. Due to the lack of preliminary symptoms and typical features, delayed diagnosis is common in clinical practice, especially in surgical systems [3]. Fragility fracture is one of the most critical complications of PMOP, leading to high disability and mortality. In China alone, the projected cost of osteoporotic fractures may reach USD 25.4 billion by 2050 [4]. Therefore, early detection is essential for alleviating the harm of PMOP.
Bone mineral density (BMD) assessed by dual-energy X-ray absorptiometry (DXA) is widely accepted as an indispensable index for defining PMOP. However, due to differences in development between developed and developing regions in medical service supply, universal access to BMD measurement seems unlikely in the short term [5]. Meanwhile, as a type of assessment method, positive imaging features always lag behind the continuous abnormality of bone metabolism, which weakens the application value of these classical examination methods in diagnosis of the early phase of bone disease [6,7]. As a direct reflection of the changes in bone homeostasis, the re-review of bone turnover markers (BTMs) in recent years seems to provide new hope for the auxiliary diagnosis of PMOP [8,9]. However, recent studies, including our own, have confirmed that there is a limited correlation between the level of serum BTMs and BMD changes [10,11]. Regretfully, few convenient and accurate biomarkers are currently available for diagnosis in the clinic.
MicroRNA (miRNA) is a type of small noncoding single-stranded RNA with 18 to 24 nucleotides [12]. As one of the epigenetic mechanisms regulating gene expression, miR-NAs mediate the posttranscriptional gene-silencing of their target genes [13]. The potential value of miRNAs as novel biomarkers for early diagnosis, treatment, and prognosis monitoring has been well verified in diseases such as cancer, obesity, and diabetes [12,14]. Although several studies have found aberrant miRNA expression in osteoporosis-induced cells and animal models [15], owing to the complexity of the pathogenesis of PMOP in humans and the imperfect public gene expression databases, the mechanisms underlying disease occurrence and progression remain to be fully elucidated. Further study on the abundance difference of miRNAs in circulating serum under pathological conditions may provide a new method to reflect the overall state of bone metabolism and the dynamic process of BMD change. Meanwhile, changes in certain miRNAs may even carry specific information about the source tissue [16] to realize the precise diagnosis and treatment of PMOP.
The primary aim of the study was to screen differentially expressed miRNAs (DEmiR-NAs) in the serum of PMOP patients and postmenopausal-without-osteoporosis (n-PMOP) controls and to validate the feasibility of using key miRNAs as biomarkers for the clinical diagnosis of disease. As a secondary aim, this study explored the expression characteristics of key miRNAs in populations with different BMDs and at different body sites. In addition, the key target gene functions and signaling pathways related to DEmiRNAs that may be involved in PMOP onset and progression were also annotated. This study highlights a novel approach in PMOP diagnosis and provides new insights into the epigenetic regulatory mechanisms of the disease.
Participants
This study surveyed 795 community-dwelling, middle-aged and elderly female participants who were recruited from June 2020 to August 2021 at the First Affiliated Hospital of Sun Yat-sen University, Guangzhou, China. The inclusion criteria were as follows: (1) age ≥ 50 years; (2) menopausal duration ≥ 1 year; and (3) signed an informed consent form before study entry. The exclusion criteria were as follows: (1) any comorbidity that could significantly affect bone metabolism, e.g., thyroid disease, diabetes, cancer, kidney disease, or ankylosing spondylitis; (2) previous treatment with anti-osteoporosis drugs or hormones (vitamin D or/and calcium supplements were allowed), e.g., estrogen or glucocorticoids; and (3) a history of tobacco smoking or alcohol dependence within the last year. Finally, a total of 83 unrelated ethnic Han Chinese women were eligible and included in the analysis.
Anthropometric Measurements
Participants wore lightweight clothing and removed their shoes before anthropometric assessments. Height and weight were measured by the corrected mechanical weight and height scale (RGZ-120, Suhong Medical Instruments Co., Changzhou, China), with an accuracy of 0.1 cm in height and 0.1 kg in weight. The average value from 3 measurements was taken for final evaluation.
Areal BMD was measured via DXA (Lunar iDXA, GE Healthcare, Chicago, IL, USA) of the lumbar spine (LS) 1-4, total hip (TH), and femoral neck (FN). All evaluations were performed by experienced diagnostic imaging physicians. The device was calibrated daily against a standard calibration phantom according to the manufacturer's instructions. Based on prior measurements, the coefficient of variation (CV) for adult measurements is 0.8% for the LS, 0.8% for the FN, and 1.4% for the TH [10].
Biochemical and Immunological Analysis
Blood samples were collected via venipuncture, with participants having fasted overnight for at least 8 h. Whole blood was left to stand at room temperature for 30 min, and serum was then collected following centrifugation at 1200× g for 10 min at 4 • C.
The analyzers were calibrated daily before the analysis of all serum samples according to the manufacturer's protocol. The routine clinical chemistry panel, including UA, ALP, calcium, and phosphorus, was detected using an AU5800 automatic biochemistry analyzer and its corresponding reagents (Beckman Coulter, Brea, CA, USA), with intra-and interassay CVs ranging from 0.5% to 4.9%. The special clinical immunology panel, including 25(OH)D, N-MID, P1NP and β-CTX, was measured using a Cobas 6000 analyzer series and its corresponding reagents (Roche, Basel, Switzerland, CH), with intra-and interassay CVs ranging from 0.6% to 4.3%.
RNA Extraction
Total RNA was extracted from 250 µL serum using TRlzol LS reagent (Invitrogen, Life Technologies, Carlsbad, CA, USA) according to the manufacturer's protocol. The RNA quantity and quality were determined using an ND-1000 Spectrophotometer (NanoDrop Technologies, Wilmington, DE, USA).
Microarray
miRNA microarray analysis was performed by a commercial service (Kangcheng Biotech Co., Shanghai, China). Briefly, miRNA expression profiling was performed using an Agilent Human miRNA Microarray system, 8 × 60 K array (Agilent Technologies, Santa Clara, CA, USA), containing probes for 2549 human miRNAs based on the miRBase database (http://www.mirbase.org, accessed on 1 December 2020, version 21.0). RNA labeling and hybridization on the Agilent miRNA microarray chips were performed with an Agilent Quick Amp Labeling Kit (Agilent part number [p/n]: 5190-0442) and Agilent Gene Expression Hybridization Kit (Agilent p/n: 5188-5242). The hybridization images were captured with an Agilent Microarray Scanner (Agilent p/n: G2565BA) and digitized using Agilent Feature Extraction (version 11.0.1.1).
The microarray data in this study have been deposited in the Gene Expression Omnibus (GEO) database (http://www.ncbi.nlm.nih.gov/geo/, accessed on 1 May 2022; Accession number: GSE201543).
qRT-PCR
To confirm the findings obtained by analyzing the miRNA profiles, qRT-PCR analysis was performed using a QuantStudio5 Real-time PCR System (Applied Biosystems, Waltham, CA, USA). cDNA was obtained from 150 ng of total RNA using M-MuLV Reverse Transcriptase (Enzymatics p/n, P7040L). The PCR amplification procedures were performed according to a previous description and repeated in triplicate [18]. The relative expression level of miRNAs was normalized to that of the internal control hsa-miR-425-5p using the 2 −∆∆Ct cycle threshold method [19]. The primer sequences for the qRT-PCR assays are listed in Supplementary Table S1.
Statistical Analysis
Statistical analyses were performed using IBM SPSS Statistics (IBM, Armonk, NY, USA; version 22.0). Data are presented as the mean ± standard deviation. Independent sample t-tests and one-way analysis of variance (ANOVA) were performed to analyze the data. Binary logistic regression analyses were performed to calculate the predicted probability of different diagnostic combinations. Receiver operating characteristic (ROC) curves were plotted to evaluate the diagnostic effects of the models. The area under the ROC curve (AUC) for different prediction models was compared using the method described by DeLong et al. [20]. To investigate associations between key miRNAs and the BMD at different body sites, Pearson's correlation and partial correlation analyses were used. Multiple linear regression models were conducted to examine the factors that influenced changes in BMD. The variance inflation factor (VIF) was calculated to assess the collinearity of independent variables, and an independent variable with VIF > 10 was considered highly collinear. Differences were considered statistically significant at p < 0.05.
Clinical Characteristics of the Participants
The characteristics of the participants in the discovery set, training set, and validation set are shown in Table 1. In the discovery set and training set, no significant differences in participant characteristics, including age, BMI, age at menopause, or menopausal duration, were observed between the PMOP patients and n-PMOP controls (p > 0.05). In the validation set, age and menopausal duration were significantly higher in the 23 PMOP patients than in the 12 n-PMOP controls (67.5 ± 8.8 vs. 58.3 ± 8.1, p = 0.005; 16.0 ± 8.5 vs. 6.3 ± 6.1, p = 0.001, respectively). The differences in BTMs and biochemical indices did not reach statistical significance between PMOP patients and n-PMOP controls in the discovery set, training set, or validation set (p > 0.05), except for calcium in the discovery set (2.10 ± 0.13 vs. 2.35 ± 0.06, p = 0.006).
Screening of Key miRNAs
The expression levels of 2549 miRNAs were measured in the discovery set. Under the criteria of adj. p < 0.05 and |Log2 Fold Change (FC)| > 1, a total of 198 DEmiRNAs were screened (Supplementary Table S2). Compared with the n-PMOP controls, 148 miR-NAs were significantly upregulated and 50 miRNAs were significantly downregulated in the PMOP patients. A volcano plot was constructed to demonstrate the profiles of the DEmiRNAs (Figure 2A), and the 50 most upregulated and downregulated miRNAs are shown in the heatmap ( Figure 2B). Among the DEmiRNAs, ten miRNAs with the highest FC in expression levels compared with the n-PMOP controls were selected as candidates for verification by qRT-PCR. The qRT-PCR results showed that the differential expression of candidate miRNAs was generally consistent with the microarray results (Supplementary Table S3).
Using the method of Delong et al. [20], the differences in the AUC among different diagnostic models were examined in the training set and validation set. The results showed that although the combination of multiple miRNAs improved the AUC to a certain extent, there was no significant difference between these AUCs in either the training set or validation set (p > 0.05).
Relative Expression Levels of Key miRNAs in Different Clinical Stages
Seventy-three participants in the training set and validation set were divided into four groups based on specific guidelines [17], and the demographics data and serum indices are shown in Supplementary Table S6. The age and menopausal duration in the severe osteoporosis group (68.7 ± 5.5 and 16.2 ± 6.5, respectively) were significantly higher than in the normal group (59.2 ± 6.3 and 5.7 ± 4.5, respectively) and osteopenia group (61.4 ± 7.8 and 10.0 ± 7.0, respectively) (p < 0.05). However, no significant difference was observed between the severe osteoporosis group and the osteoporosis group in terms of demographics (p > 0.05). ALP levels (108.2 ± 36.3) in the serum of the severe osteoporosis group were the highest among the four groups (p < 0.05). N-MID, ALP, and phosphorus were the only three serum indices that showed significant differences among the four groups (p < 0.05).
The relative expression levels of the key miRNAs in serum are shown in Figure 4. hsa-miR-144-5p, hsa-miR-506-3p, and hsa-miR-6851-3p were highly expressed in the osteoporosis group, and the differences were significant compared with the normal group and osteopenia group (p < 0.05). In addition, hsa-miR-144-5p and hsa-miR-8068 were expressed at lower levels in the osteopenia group than in the severe osteoporosis group (p < 0.05). Differences in key miRNA serum levels among subgroups may partly reflect the unique bone metabolic patterns in various stages of disease.
Correlations between Key miRNAs and BMD
According to the results of Pearson correlation analysis, hsa-miR-144-5p was significantly negatively correlated with FN BMD (p < 0.05), but there was no significant correlation with the BMD at other body sites (p > 0.05); hsa-miR-506-3p and hsa-miR-8068 were significantly associated with LS 1-4 BMD (p < 0.05), while hsa-miR-6851-3p was significantly correlated with BMD at all three body sites (p < 0.05). Additionally, age and menopausal duration were used as covariates for partial correlation analysis. Apart from hsa-miR-8068, the other candidate key miRNAs had different degrees of negative correlations with BMD at LS 1-4, TH, and FN (p < 0.05), and the correlations were further enhanced (Table 3). To further evaluate the influencing factors of the key miRNAs on BMD at different body sites, multivariable linear regression was performed using the key miRNAs as independent variables and BMD as a dependent variable. Meanwhile, variables that were significant in univariate analyses were adopted as covariates. No multicollinearity was detected among the independent variables in the regression models. Table 4 shows the results from the linear regression models. In Model 1, only hsa-miR-6851-3p was found to be significantly negatively associated with LS 1-4 BMD (β = −0.851, p = 0.007). When controlling for age, menopausal duration, N-MID, ALP, and phosphorus as covariates in Model 2, hsa-miR-6851-3p remained a strong predictor of LS 1-4 BMD (β = −0.645, p = 0.026), while hsa-miR-144-5p was determined to be the lone predictor to have significant predictive power for BMD at different body sites (LS 1-4: β = −0.265, p = 0.022; TH: β = −0.301, p = 0.005; and FN: β = −0.324, p = 0.003, respectively). No significant correlation was found between BMD and hsa-miR-506-3p or hsa-miR-8068 (p > 0.05).
Target Genes and Pathways Correlated with DEmiRNAs
On the basis of the 198 DEmiRNAs, 1945 target genes were identified using the TargetScan, miRTarBase, and miRDB databases. The top 10 miRNAs and their target genes are shown in Supplementary Table S7. In total, 1120 GO terms and 85 KEGG pathways were enriched in target genes according to the criteria adj. p < 0.05 (Supplementary Tables S8 and S9). The top 10 enriched terms of GO and KEGG pathways are shown in Figure 5A,B. The 1945 target genes were imported into the STRING database to construct a PPI network. After excluding the isolated nodes, the final PPI network is composed of 1052 nodes and 7477 edges, as is shown in Supplementary Figure S2. The top 20 hub genes were sorted by degree (refers to the number of gene connections within the network) in descending order (Supplementary Table S10). Ninety-four target genes for the five candidate key miRNAs were also predicted. As shown in Figure 5C, hsa-miR-340-5p and hsa-miR-144-5p cooperatively regulate the target gene CREBRF. SYNPO2 is the only predicted target gene of hsa-miR-144-5p. The function of target genes was significantly enriched in 17 GO terms (Biological Process [BP]:Cellular Component [CC]:Molecular Function [MF] = 8:4:5) and a KEGG signaling pathway ("neurotrophin signaling pathway"). For the BP classification, the top three enrichment terms were "regulation of release of sequestered calcium ion into cytosol by sarcoplasmic reticulum", "regulation of ryanodine-sensitive calcium-release channel activity", and "regulation of heart rate". Membrane, focal adhesion, and cytosol were the most highly enriched CC terms. In the MF category, "protein binding", "Smad binding", and "poly(A) RNA binding" were significantly enriched ( Figure 5D). For visualization, the PPI network of target genes is presented in Figure 5E. YY1, VIM, and YWHAE were the top three hub genes, with higher interaction levels ( Table 5).
Gene
Degree Gene Degree
Discussion
Since the first discovery of miRNAs in Caenorhabditis elegans in 1993, these molecules have been widely confirmed to exist in more than 12 types of mammalian body fluids, including serum [12,21]. The high degree of conservation [22], detectability [23], and specific spatiotemporal expression [14] of serum miRNAs make them promising biomarkers for liquid biopsy. As important regulators of gene expression, an increasing number of studies have confirmed that miRNAs are involved in regulating the pathological progression of PMOP [24][25][26]. However, because of the absence of clinical validation, none of them has been widely recommended for clinical diagnostic purposes. In this study, four key miRNAs were screened and validated in independent populations. The results revealed that hsa-miR-144-5p, hsa-miR-506-3p, hsa-miR-8068, and hsa-miR-6851-3p are potential independent biomarkers for distinguishing PMOP patients from n-PMOP controls. This finding provides additional information for PMOP diagnosis, independent of radiological findings and BTMs.
Several small clinical studies have preliminarily explored the potential of abnormally expressed circulating miRNAs as diagnostic biomarkers of PMOP [27][28][29]. For example, miR-133a, a stimulator of osteoclastogenesis, is useful for the detection of PMOP [29]. However, the results of previous studies are inconsistent, possibly due to the different populations studied. Moreover, the diagnostic accuracy and clinical usability of certain miRNAs for disease in certain clinical settings are limited. A study by Mandourah et al. reported that hsa-miR-122-5p and hsa-miR-4516 were downregulated in blood samples and could serve as potential diagnostic markers of osteoporosis, but the diagnostic values are not yet sufficient (AUC = 0.752) [28]. In our study, four key miRNAs showed high diagnostic value, with AUCs greater than 0.7 in both the training set and validation set, and the AUCs using the combined miRNAs for diagnosis were higher than 0.9.
Currently, the functions of hsa-miR-144-5p and hsa-miR-506-3p have been reported mainly in the cancer field. High plasma levels of hsa-miR-144-5p were shown to be associated with renal cell carcinoma [30], non-small-cell lung cancer [31] or glioblastoma [32]. Recent research by Zhang et al. showed that miR-144-5p can reduce bone repair and regeneration in type 2 diabetes by suppressing the expression of Smad1 [33]. In this study, linear regression analysis revealed that hsa-miR-144-5p was independently associated with BMD changes in multiple body sites, and the association between its expression pattern and PMOP progression is promising for further longitudinal research. Overexpression of hsa-miR-506-3p was found to inhibit the proliferation, migration, and invasion of cancer cells in osteosarcoma [34], prostate cancer [35], and hepatocellular carcinoma [36]. Thus far, the function and serum expression profiles of hsa-miR-8068 and hsa-miR-6851-3p have not been extensively investigated. The results of the present study demonstrate for the first time the additional values of the above miRNAs in PMOP.
The predicted expression profiles of the target genes confirmed the close association between DEmiRNAs and bone metabolism. As reported by Jeong et al., YY1 significantly inhibited Runx2-mediated transcriptional activity of osteocalcin (OCN) and ALP promoters, and knockout of this gene enhanced the osteoblast differentiation induced by BMP2 and Runx2 [37]. At the same time, YY1 can modulate the transcriptional activity of Smad, thereby regulating cell differentiation induced by the transforming growth factor superfamily signaling pathway [38]. VIM is a type III intermediate filament protein that is expressed in mesenchymal cells [39]. Overexpression of VIM in osteoblasts inhibits osteoblast differentiation, as shown by reduced ALP activity, delayed mineralization, and reduced expression of osteoblast markers [40]. This effect may be mediated by VIM competitively binding ATF4, which is required for OCN transcription and osteoblast differentiation [41]. YWHAE belongs to the 14-3-3 protein family and is involved in the transduction of signaling pathways by binding to phosphoserine-containing proteins [42]. A study on YWHAE in exosomes released from an osteoblast/osteocyte coculture system revealed that YWHAE had a positive response to mechanical stress [43]. Rivero et al. loaded 14-3-3ε protein into a scaffold, and positive stimulation of osteogenicity was observed [44]. These reports regarding the functions of target genes predicted by DEmiRNA profiles support our finding that elevated levels of differentially expressed serum miRNAs are correlated with PMOP.
Despite its strengths as discussed above, this research had certain limitations. First, we used strict inclusion criteria to obtain relatively homogenous populations. However, a power calculation was not performed a priori, and the small sample size may have resulted in decreased statistical power. Second, despite accounting for many important confounders, we cannot exclude residual confounding by unmeasured or unknown confounders. Thus, these findings need to be validated in subsequent large sample studies. In addition, further functional validation in vitro and in vivo is required to confirm the associations between these molecules and disease.
Conclusions
This study presents new clinical evidence regarding the deregulated expression of miRNAs in the serum of PMOP patients. Our results indicate that hsa-miR-144-5p, hsa-miR-506-3p, hsa-miR-8068, and hsa-miR-6851-3p target a variety of bone-metabolism-related genes and pathways and are potential independent biomarkers for clinical diagnosis of the disease, outperforming traditional BTMs. Among them, hsa-miR-144-5p is the only key miRNA that can simultaneously predict changes in BMD in LS 1-4, TH, and FN. These findings provide not only a new method for clinicians to evaluate the changes in BMD in postmenopausal women under limited conditions but also a meaningful inspiration for in-depth study of the epigenetic regulation mechanisms underlying PMOP onset and progression. Informed Consent Statement: Informed consent was obtained from all subjects involved in the study.
Data Availability Statement:
The datasets supporting the conclusions of this article are included within the paper and its Supplementary Materials. All other datasets used and analyzed during the study are available from the corresponding author on reasonable request. | 5,294.8 | 2022-11-01T00:00:00.000 | [
"Biology",
"Medicine"
] |
Research on Road Scene Understanding of Autonomous Vehicles Based on Multi-Task Learning
Road scene understanding is crucial to the safe driving of autonomous vehicles. Comprehensive road scene understanding requires a visual perception system to deal with a large number of tasks at the same time, which needs a perception model with a small size, fast speed, and high accuracy. As multi-task learning has evident advantages in performance and computational resources, in this paper, a multi-task model YOLO-Object, Drivable Area, and Lane Line Detection (YOLO-ODL) based on hard parameter sharing is proposed to realize joint and efficient detection of traffic objects, drivable areas, and lane lines. In order to balance tasks of YOLO-ODL, a weight balancing strategy is introduced so that the weight parameters of the model can be automatically adjusted during training, and a Mosaic migration optimization scheme is adopted to improve the evaluation indicators of the model. Our YOLO-ODL model performs well on the challenging BDD100K dataset, achieving the state of the art in terms of accuracy and computational efficiency.
Introduction
The compositions of autonomous driving vehicles can be divided into three modules: environmental perception, decision planning, and vehicle control. Environmental perception is the most fundamental part to realize autonomous driving and is one of the critical technologies in intelligent vehicles [1]. The performance of perception will determine whether the autonomous vehicle can adapt to the complex and changeable traffic environment. The research progress of computer vision shows that visual perception will play a decisive role in the development of autonomous driving [2]. Moreover, vision sensors have the advantages of mature technologies, low prices, and comprehensive detection [3,4].
Effective and rapid detection of traffic objects in various environments will ensure the safe driving of autonomous vehicles, but the detection performance is severely limited by road scenes, lighting, weather, and other factors. The development of big data, computing power, and algorithms has continuously improved the accuracy of deep learning.
Great breakthroughs have been made in the automatic driving industry, making the above detection problems expected to be solved. At present, creating deeper learning networks as deep as possible is the main trend in current research [5], and, while significant progress has been made in that direction, the demands on computational power are becoming increasingly demanding. Runtime is becoming very important when it comes to actually deploying applications.
In view of this, aiming at the problems of redundant calculation, slow speed, and low accuracy in the existing perception models, we propose a multi-task model YOLO-Object, Drivable Area, and Lane Line Detection (YOLO-ODL) model based on hard parameter sharing, which realizes the joint and efficient detection of objects, drivable areas, and At present, the YOLOv5 object detection model has the best tradeoff between detection accuracy and speed [6]. Based on the YOLOv5s model [7], we built the YOLO-ODL multi-task model. YOLO-ODL has a good balance of detection speed and accuracy, achieving state-of-the-art performance on the BDD100K dataset [8], reaching 94 FPS on an RTX 3090 GPU. The detection speed can reach 477 FPS using TensorRT.
The main contributions of this work are as follows: (I) On the basis of the YOLOv5s model, a multi-task model YOLO-ODL based on hard parameter sharing is built to realize joint and efficient detection of traffic objects, drivable areas, and lane lines. (II) The performance of traffic object detection task is improved by adding shallow highresolution features and changing the size of the output's feature map. (III) In order to further improve the performance of YOLO-ODL, the weight balance strategy and Mosaic migration optimization scheme are introduced to improve the evaluation indicators of the multi-task model effectively.
The structure of this paper is organized as follows. Section 2 analyzes related work. Section 3 presents the YOLO-ODL multi-task model. Section 4 presents experimental verification results. Section 5 presents conclusions.
Related Work
In this section, we review related solutions for the above three tasks of object detection, drivable area detection, and lane line detection, respectively, and then introduce some related multi-task learning work.
Object Detection
Transformer-based object detection methods have had dominant performances in recent years. Zhu et al. [9] applied the transformer to YOLO and achieved be er detection results. Anchor-based methods are still the mainstream of object detection at present [10]. Their core idea is to introduce anchor boxes, which can be considered as pre-defined proposals, as a priori for bounding boxes, which can be divided into one-stage object detection methods and two-stage object detection methods. As it is necessary to extract object regions from a set of object proposals and then classify them, the two-stage method is less efficient than the one-stage method [11]. The one-stage method uses the one-stage network to directly output the location and category of objects so that it has evident advantages in training and reasoning time. Liu et al. [12] presented a method of detecting objects in images using a single depth neural network SSD, which could improve the detection speed and achieve the accuracy of two-stage detectors at the same time. Lin et al. [13] designed and trained a simple dense detector called RetinaNet that introduced a new At present, the YOLOv5 object detection model has the best tradeoff between detection accuracy and speed [6]. Based on the YOLOv5s model [7], we built the YOLO-ODL multitask model. YOLO-ODL has a good balance of detection speed and accuracy, achieving state-of-the-art performance on the BDD100K dataset [8], reaching 94 FPS on an RTX 3090 GPU. The detection speed can reach 477 FPS using TensorRT.
The main contributions of this work are as follows: (I) On the basis of the YOLOv5s model, a multi-task model YOLO-ODL based on hard parameter sharing is built to realize joint and efficient detection of traffic objects, drivable areas, and lane lines. (II) The performance of traffic object detection task is improved by adding shallow highresolution features and changing the size of the output's feature map. (III) In order to further improve the performance of YOLO-ODL, the weight balance strategy and Mosaic migration optimization scheme are introduced to improve the evaluation indicators of the multi-task model effectively.
The structure of this paper is organized as follows. Section 2 analyzes related work. Section 3 presents the YOLO-ODL multi-task model. Section 4 presents experimental verification results. Section 5 presents conclusions.
Related Work
In this section, we review related solutions for the above three tasks of object detection, drivable area detection, and lane line detection, respectively, and then introduce some related multi-task learning work.
Object Detection
Transformer-based object detection methods have had dominant performances in recent years. Zhu et al. [9] applied the transformer to YOLO and achieved better detection results. Anchor-based methods are still the mainstream of object detection at present [10]. Their core idea is to introduce anchor boxes, which can be considered as pre-defined proposals, as a priori for bounding boxes, which can be divided into one-stage object detection methods and two-stage object detection methods. As it is necessary to extract object regions from a set of object proposals and then classify them, the two-stage method is less efficient than the one-stage method [11]. The one-stage method uses the one-stage network to directly output the location and category of objects so that it has evident advantages in training and reasoning time. Liu et al. [12] presented a method of detecting objects in images using a single depth neural network SSD, which could improve the detection speed and achieve the accuracy of two-stage detectors at the same time. Lin et al. [13] designed and trained a simple dense detector called RetinaNet that introduced a new loss function Focal Loss, which effectively solved the problem of imbalanced proportions of positive and negative samples in the training process. The subsequent versions of YOLO [7,14,15] a series of improvement measures to further improve the object detection accuracy on the premise of ensuring the detection speed.
The actual running time of the object detection algorithm is very important for deploying the application online, so it is necessary to further balance the detection speed and accuracy. In addition, there are a large number of small objects in the autonomous driving environment, but, due to the low-resolution of small objects and the lack of sufficient appearance and structure information, small-object detection is more challenging than ordinary-object detection.
Drivable Area Detection
With the rapid development of deep learning technologies, many effective semantic segmentation methods have been proposed and applied to drivable area detection [16]. Long et al. [17] first proposed FCN, which is an end-to-end semantic segmentation network for opening a precedent for using convolutional neural networks to deal with semantic segmentation problems. Badrinarayanan et al. [18] proposed an encoder-decoder semantic segmentation network named SegNet, which is currently widely used. Zhao et al. [19] established PSPNet, a network that extends pixel-wise features to global pyramid pooling features, thereby combining context information to improve detection performance. Unlike the above-mentioned networks, Tian et al. [20] proposed decoder-adopted, data-dependent upsampling (DUpsampling), which can recover the resolution of feature maps from the low-resolution output of the network. Takikawa et al. [21] proposed a new structure GSCNN, which uses a new structure composed of shape branches and rule branches to focus on boundary information and improves the segmentation ability of small objects.
Based on deep learning, semantic segmentation effectively improves detection accuracy. However, due to the increasing number of deep learning network layers, the model becomes more complex, which leads to the low efficiency of segmentation and is difficult to apply to autonomous driving scenarios.
Lane Line Detection
In order to significantly improve the accuracy and robustness of lane line detection, many lane line detection methods based on deep learning have been proposed. Lane line detection methods can be roughly divided into two categories; one is based on classification, the other is based on semantic segmentation. The classification-based lane detection method reduces the size and computation of the model but suffers from a loss in accuracy and cannot detect scenes with many lane lines well [22]. The method based on semantic segmentation classifies each pixel into lane or background. Pan et al. [23] proposed Spatial CNN (SCNN) that replaced the traditional layer-by-layer convolutions with slice-by-slice convolutions, which enabled message passing between pixels across rows and columns in a layer, thereby improving the segmentation ability for long continuous shape structures such as lane lines, poles, and walls. Hou et al. [24] presented a novel knowledge distillation approach named the Self Attention Distillation (SAD) mechanism, which was incorporated into a neural network to obtain significant improvements without any additional supervision or labels. Zheng et al. [25] presented a novel module named REcurrent Feature-Shift Aggregator (RESA) to collect feature map information so that the network could transfer information more directly and efficiently.
Lane line detection based on semantic segmentation can effectively increase detection accuracy. However, with the increase in detection accuracy, the network becomes more complex, and each pixel needs to be classified, so the detection speed needs to be improved.
Multi-Task Learning
Multi-task networks usually adopt the scheme of hard parameter sharing, which consists of a shared encoder and several feature task decoders. The proposed multi-task network [26] has the advantages of small size, fast speed, high accuracy, and can be used for positioning, making it highly suitable for online deployment of autonomous vehicles. Teichmann et al. [2] proposed an efficient and effective feedforward architecture called MultiNet, which combined classification, detection, and semantic segmentation; the approach shares a common encoder and has three branches on which specific tasks are built with multiple convolutional layers. DLT-Net, proposed by Qian et al. [1], and YOLOP, proposed by Wu et al. [27], simultaneously deal with the problems of traffic objects, drivable areas, and lane line detections in one framework. A multi-task network is helpful to improve model generalization and reduce computing costs. However, multi-task learning may also reduce the accuracy of the model due to the need to balance multiple tasks [28].
Models based on multi-task learning need to learn knowledge from different tasks and are highly dependent on the shared parameters of multi-task models. However, the multitask architectures proposed by most of the current works lack balancing the relationships among the various tasks. Multi-task models are usually difficult to train to achieve the best effect, and unbalanced learning will reduce the performance of a multi-task model, so it is necessary to adopt meaningful feature representation and appropriate balanced learning styles.
Proposed Method
Considering the advantages of the multi-task model with hard parameter sharing, such as simple structure, high operating efficiency, and low over-fitting risk, a multi-task model YOLO-ODL based on hard parameter sharing is proposed to achieve joint and efficient detection of objects, drivable areas, and lane lines. The structure of the YOLO-ODL model is shown in Figure 2, which consists of one shared encoder (Shared Encoder) and three decoders (Detection Encoder, Drivable Area Encoder, and Lane Line Encoder) to solve specific tasks.
Multi-Task Learning
Multi-task networks usually adopt the scheme of hard parameter sharing, which consists of a shared encoder and several feature task decoders. The proposed multi-task network [26] has the advantages of small size, fast speed, high accuracy, and can be used for positioning, making it highly suitable for online deployment of autonomous vehicles. Teichmann et al. [2] proposed an efficient and effective feedforward architecture called MultiNet, which combined classification, detection, and semantic segmentation; the approach shares a common encoder and has three branches on which specific tasks are built with multiple convolutional layers. DLT-Net, proposed by Qian et al. [1], and YOLOP, proposed by Wu et al. [27], simultaneously deal with the problems of traffic objects, drivable areas, and lane line detections in one framework. A multi-task network is helpful to improve model generalization and reduce computing costs. However, multi-task learning may also reduce the accuracy of the model due to the need to balance multiple tasks [28].
Models based on multi-task learning need to learn knowledge from different tasks and are highly dependent on the shared parameters of multi-task models. However, the multi-task architectures proposed by most of the current works lack balancing the relationships among the various tasks. Multi-task models are usually difficult to train to achieve the best effect, and unbalanced learning will reduce the performance of a multitask model, so it is necessary to adopt meaningful feature representation and appropriate balanced learning styles.
Proposed Method
Considering the advantages of the multi-task model with hard parameter sharing, such as simple structure, high operating efficiency, and low over-fi ing risk, a multi-task model YOLO-ODL based on hard parameter sharing is proposed to achieve joint and efficient detection of objects, drivable areas, and lane lines. The structure of the YOLO-ODL model is shown in Figure 2, which consists of one shared encoder (Shared Encoder) and three decoders (Detection Encoder, Drivable Area Encoder, and Lane Line Encoder) to solve specific tasks.
Shared Encoder
Shared Encoder shared neural network parameters and adopted Backbone and FPN structures in the YOLOv5s object detection model. Backbone extracted common image features from scenes, and FPN integrated image features of different scales (i.e., the information required to detect objects, driving areas, and lane lines). To enhance the feature
Shared Encoder
Shared Encoder shared neural network parameters and adopted Backbone and FPN structures in the YOLOv5s object detection model. Backbone extracted common image features from scenes, and FPN integrated image features of different scales (i.e., the information required to detect objects, driving areas, and lane lines). To enhance the feature extraction capability of the Shared Encoder, a shallow high-resolution feature with a size of 160 × 96 was added to the FPN structure. The input image size of the model was 640 × 384 × 3, and the Shared Encoder generated feature maps with three sizes of 160 × 96 × 128, 80 × 48 × 64, and 40 × 24 × 128 from bottom up.
Detection Decoder
Detection Decoder included PANet and Detection Head structures in YOLOv5s object detection model, which were used to decode object detection tasks. PANet further integrated features of different scales, and Detection Head adopted convolutional layer to adjust the number of channels, with 1 × 1 kernel size and stride equal to 1. A lot of small objects, such as traffic lights, traffic signs, pedestrians, and vehicles in the distance, exist in autonomous driving scenarios. Shallow high-resolution features are very important to detect small objects [29], so shallow high-resolution features with sizes of 160 × 96 were added to replace the initial deep low-resolution features with sizes of 20 × 12, and, finally, feature maps with three sizes of 160 × 96 × 18, 80 × 48 × 18, and 40 × 24 × 18 were generated. As each grid was responsible for three anchor boxes, there were a total of 61,440 prediction outputs, and each prediction output included four parameters related to the position of the prediction box, one confidence parameter, and one vehicle category parameter, so the output feature map had 18 channels.
Drivable Area Decoder and Lane Line Decoder
Inspired by the DLT-Net [1] and YOLOP [27] models, we adopted semantic segmentation to realize the detection task of drivable areas and lane lines and used the same decoding structure. The purpose of decoding was to transform the intermediate feature map back to the resolution size of the input image while reducing the number of feature map channels, i.e., transform a low-resolution feature map with a size of 160 × 96 × 128 back to a high-resolution feature map with a size of 640 × 384 × 2, two channels corresponded to the number of classifications, and image features are further extracted to generate denser feature maps, and, finally, the semantic probability output of drivable area and lane line segmentation was generated through the Sigmoid layer.
Loss Function
In the training process of a multi-task model, all tasks start to learn at the same time, so it is necessary to correlate the losses of multiple tasks. We defined the total loss as a weighted sum of losses for object detection, drivable area detection, and lane line detection tasks, so as to make the loss scales of each task closer, where object detection included 3 subtasks concerning bounding box, category, and confidence. The total loss was expressed as L total = ∑ i∈ box, cls, obj, drivable, lane where θ i is the neural network parameters, and w i and L i are the loss weight and loss function of specific tasks, respectively. Remark 1: (1) uses a weight balancing scheme to achieve dynamic adjustment of the weight parameters so as to balance the relationship between individual tasks. The local minima of different tasks in multi-task learning are in different locations, and, by interacting with each other, they can help to escape the local minima, and this paper uses a weight balancing scheme to achieve dynamic adjustment of the weight parameters to minimize the overall loss.
Bounding box prediction is about the judgment of regression parameters, which is a regression problem. From the traditional Smooth L1 to CIoU (Complete IoU) regression loss function, the prediction speed and accuracy of the bounding box have been greatly improved. Therefore, the CIoU loss function is adopted in this paper, and the CIoU loss is calculated by using prediction frames and real frames. The important geometric factors considered by the CIoU Loss function include the overlapping area, center distance, and aspect ratio of the prediction frame and the real frame. The formula for the loss of the bounding box L box is as follows where the coordinates of the center point and width and height of the predicted frame are (x,ŷ) and ŵ,ĥ , respectively, and the coordinates of the center point and width and height of the real frame are (x, y) and (w, h), respectively. ρ represents the Euclidean distance between the center point of the predicted frame and the real frame; c represents the diagonal length of the smallest outer enclosing rectangle of the predicted frame and the real frame; and v measures the similarity of the aspect ratio between the predicted frame and the real frame.
Confidence prediction is to judge whether the boundary box contains the target, and category prediction is to judge the target category in the boundary box, both of which belong to classification problems. In this paper, the Binary Cross Entropy (BCE) loss function was used to calculate the classification loss, and the BCE loss function required the input data to be in the range of [0, 1], so the Sigmoid function needed to be used to standardize the input data before input. Assuming that the input image was divided into S × S grids, and each grid was responsible for B prior boxes, the calculation formula of confidence loss L obj and classification loss L cls was as follows ∑ nc∈classes P(i, j, nc) · logP(i, j, nc) +(1 − P(i, j, nc)) · log 1 −P(i, j, nc) where C represents the confidence truth label;Ĉ represents the confidence level of the prediction; obj represents the prediction box corresponds to the target object; noobj represents the prediction box does not correspond to the target object; I obj ij represents that if the jth prior box of the i-th grid is responsible for the target, its value is 1, otherwise it is 0; num(I obj ij ) represents the number of features corresponding to a target; I noobj ij represents that if the j-th prior box of the i-th grid is not responsible for the target, its value is 1, otherwise it is 0; num(I noobj ij ) represents the number of features corresponding to no target; nc represents the number of target detection categories number; P represents the category true label;P represents the predicted category; and the left and right terms in the L obj formula are the confidence losses for positive and negative samples, respectively.
Semantic segmentation is the classification of each pixel in the image, which belongs to the classification problem. Binary cross entropy (BCE) loss function was used to calculate, and Sigmoid function was also used to convert the input data of BCE loss function to [0, 1], where the loss of the whole graph was the average of the loss of each pixel. Semantic segmentation loss was composed of driving area detection loss L drivable and lane line detection loss L lane , and the same loss function was used to calculate segmentation loss; the calculation formula is as follows: S(m, n) · logŜ(m, n) +(1 − S(m, n)) · log 1 −Ŝ(m, n) where W and H represent the width and height of the final output feature map of the segmentation model; S represents the semantic segmentation true label; andŜ represents the predicted semantic information.
The bounding box prediction belongs to the regression problem, L box adopts the CIoU loss function, and other tasks belong to the classification problem, using the BCE loss function.
Due to the uncertainty and complexity of different tasks, the multi-task model often appears as a dominant task in the training stage, resulting in the phenomenon of unbalanced training. The performance of the multitask model depends on the weight selection of loss to a large extent, whereas manually adjusting the weight of loss is a time-consuming and laborious process. Cipolla et al. [30] introduced homoscedastic uncertainty to balance multiple tasks and added a learnable noise parameter σ to the loss of each task so that the multi-task network could automatically adjust the weight parameters during training, so as to balance various tasks. Compared with the loss-weighted summation method, this balancing method not only eliminated the need to manually adjust the weight parameters but also improved the joint training performance of the multi-task network. The multi-task loss we finally adopted was shown as
Experimental Setup
(1) Dataset: The proposed model dataset was trained and tested on the BDD100K [8], a challenging dataset that included a diverse set of driving data under various cities, weather conditions, times, and scene types. The BDD100K dataset also came with a rich set of annotations, including object bounding boxes, drivable areas, and lane markings, with a total of 100 K images with resolution sizes of 1280 × 720, including training (70 K), validation (10 K), and testing (20 K) sets. As the test set was not labeled, we tested the proposed model on the verification set. (2) Metrics: The object detection performance was evaluated by Recall and Mean Average Precision (mAP). mAP50 (mean Average Precision) represented the average precision value of the Intersection over Union (IoU) threshold of 0.5. The drivable area detection performance was evaluated by MIoU [31]. Lane line detection performance was evaluated by pixel accuracy and IoU of lanes [24]. (3) Implementation Details: The experimental environment used the Ubuntu 18.04 operating system; a GeForce RTX 3090 graphics card for computing, with 24 GB video memory size; the CPU configuration was an Intel i7-11700K @ 3.60 GHz; the CUDA version was 11.1; the PyTorch version was 1.8.0; and the Python version was 3.8. All experiments were carried out in the same experimental environment with the same training parameters. The initial learning rate was set to 0.001, the weight decay was 0.0005, the momentum was 0.937, and the Adam optimizer was used for optimization training. We adjusted the input size of the model to 640 × 384 to speed up the model and normalize the input image at the same time.
Since we only used BDD100K data set for training, the data were augmented in order to avoid overfitting and improve the generalization ability of the model. The data augmentation adopted included rotating, scaling, translation, color space augmentation, and left-right flipping. The data augmentation process is shown in Figure 3. training. We adjusted the input size of the model to 640 × 384 to speed up the model and normalize the input image at the same time.
Since we only used BDD100K data set for training, the data were augmented in order to avoid overfi ing and improve the generalization ability of the model. The data augmentation adopted included rotating, scaling, translation, color space augmentation, and left-right flipping. The data augmentation process is shown in Figure 3.
Experimental Evaluation
The object detection, drivable area, and lane line branches of the decoder were trained separately, and then the three decoder branches above were trained jointly and compared with various public models. The experiments showed that the single-task model proposed was competitive, and the jointly trained multi-task model reached the state-of-the-art level.
(1) Object Detection Branch: In order to compare with other public models, we filtered out the vehicle class annotations from the BDD100K dataset and trained the singletask model YOLO-Object Detection (YOLO-O) for object detection.
Remark 2: The proposed YOLO-O was improved on the basis of the YOLOv5s model, adding shallow high-resolution features to effectively improve the ability of the target detection, and the detection accuracy was significantly higher than YOLOv5s. Table 1 shows the comparison of object detection performance. It can be seen that YOLO-O had a higher recall value and was 24.6% and 3% higher than Faster R-CNN and YOLOv5s in terms of mAP50, respectively. It was further proved that the improved object detection model could effectively improve the performance of object detection. (2) Drivable Area Branch: We uniformly labeled the drivable areas (area/drivable) and the alternative driving areas (area/alternative) in the BDD100K dataset as drivable areas and trained the single-task model YOLO-Drivable Area Detection (YOLO-D) for drivable area detection. Table 2 shows the comparison of drivable area detection performance. It can be seen that YOLO-D was 23.7% and 2.8% higher than ERFNet and PSPNet in terms of MIoU, respectively.
Experimental Evaluation
The object detection, drivable area, and lane line branches of the decoder were trained separately, and then the three decoder branches above were trained jointly and compared with various public models. The experiments showed that the single-task model proposed was competitive, and the jointly trained multi-task model reached the state-of-the-art level.
(1) Object Detection Branch: In order to compare with other public models, we filtered out the vehicle class annotations from the BDD100K dataset and trained the single-task model YOLO-Object Detection (YOLO-O) for object detection.
Remark 2: The proposed YOLO-O was improved on the basis of the YOLOv5s model, adding shallow high-resolution features to effectively improve the ability of the target detection, and the detection accuracy was significantly higher than YOLOv5s. Table 1 shows the comparison of object detection performance. It can be seen that YOLO-O had a higher recall value and was 24.6% and 3% higher than Faster R-CNN and YOLOv5s in terms of mAP50, respectively. It was further proved that the improved object detection model could effectively improve the performance of object detection. (2) Drivable Area Branch: We uniformly labeled the drivable areas (area/drivable) and the alternative driving areas (area/alternative) in the BDD100K dataset as drivable areas and trained the single-task model YOLO-Drivable Area Detection (YOLO-D) for drivable area detection. Table 2 shows the comparison of drivable area detection performance. It can be seen that YOLO-D was 23.7% and 2.8% higher than ERFNet and PSPNet in terms of MIoU, respectively.
Model MIoU
ERFNet [1] 68.7% PSPNet [27] 89.6% YOLO-D 92.4% (3) Lane Line Branch: As the lanes of the BDD100K dataset were labelled by two lines, we used the ENet-SAD [24] data enhancement (see Figure 3 for details) to effectively avoid over fitting and improve the generalization ability of the model, which had advantages in the detection of driving areas and lane lines. Table 3 shows the comparison of lane line detection performance, where CGAN-L adopted a conditional generative adversarial network to detect lane lines. Although YOLO-L was 1.9% lower than CGAN-L in terms of IoU, it was much higher than the other four lane line detection models in terms of accuracy. Table 3. Lane Line Detection Performance Comparison.
Model Accuracy IoU
SCNN [27] 35.8% 15.8% ENet-SAD [24] 36.6% 16.0% CGAN-L [32] 57.2% 30.0% SALMNet [33] 58.3% 25.1% YOLO-L 80.0% 28.1% (4) YOLO-ODL: We trained the YOLO-ODL multi-task model to realize the joint detection of objects, drivable areas, and lane lines. Table 4 shows the comparison of YOLO-ODL performance indicators. It can be seen that the runtime of the YOLO-ODL model had evident advantages over the total runtimes of the three single-task models, but it was a little worse in terms of evaluation indicators. Therefore, we adopted the weight balancing scheme of Formula (2) to train the YOLO-ODL model so that the model could automatically balance the multi-task weights. Mosaic data augmentation effectively improved the accuracy of object detection [15], so the migration optimization scheme shown in Figure 4 was adopted. Mosaic data augmentation was added, and then the knowledge transfer learned from the YOLO-O object detection model was applied to the multi-task model. Remark 3: The proposed YOLO-D and YOLO-L models com the improved YOLOv5s model to expand the semantic segmenta data enhancement (see Figure 3 for details) to effectively avoid the generalization ability of the model, which had advantages in areas and lane lines. Table 3 shows the comparison of lane line detection perfor adopted a conditional generative adversarial network to dete YOLO-L was 1.9% lower than CGAN-L in terms of IoU, it was mu four lane line detection models in terms of accuracy.
Model
Accuracy SCNN [27] 35.8% ENet-SAD [24] 36.6% CGAN-L [32] 57.2% SALMNet [33] 58.3% YOLO-L 80.0% (4) YOLO-ODL: We trained the YOLO-ODL multi-task model tion of objects, drivable areas, and lane lines. Table 4 shows th ODL performance indicators. It can be seen that the run model had evident advantages over the total runtimes of th els, but it was a li le worse in terms of evaluation indicators the weight balancing scheme of Formula (2) to train the YO the model could automatically balance the multi-task weigh tation effectively improved the accuracy of object detection [ timization scheme shown in Figure 4 was adopted. Mosaic added, and then the knowledge transfer learned from the Y model was applied to the multi-task model. As shown in Table 4, YOLO-ODL (+ weight balance) improved the performances of the three tasks, which could better balance each task. YOLO-ODL (+ weight balance and migration optimization) further improved the evaluation indicators of the model, especially in object detection, thus proving the effectiveness of weight balance and migration optimization. In general, the optimized YOLO-ODL had advantages over the three single-task models.
In view of the variability of weather and the complexities of road scenes, the YOLO-ODL model was tested under multi-weather and multi-scenario road conditions, respectively, to further embody the robustness of the acceleration model. As shown in Figures 5 and 6, YOLO-ODL had strong robustness in multi-weather and multi-scenario road conditions and could well detect traffic objects, drivable areas, and lane lines.
detection, the proposed YOLO-ODL was superior to other multi-task models.
As shown in Figures 7-9, we compared the detection effects of YOLO-ODL and YOLOP on three tasks. Object detection results are shown in Figure 7. YOLOP was prone to false detections and missed detections, mistakenly detecting complex house backgrounds as vehicles, missing the detection of vehicles in the distance. Figure 8 shows drivable area detection results. YOLOP had incomplete detection and many false detection areas, whereas the drivable area detected by YOLO-ODL was more accurate and complete. Lane line detection results are shown in Figure 9. The lane lines detected by YOLO-ODL were more accurate and continuous, but the lane lines detected by YOLOP were discontinuous, and there were more noise areas. The comparison of the multi-task model detection performance is shown in Table 5. As the experimental test platforms were different, the detection speeds of MultiNet and DLT-Net were much lower than that of YOLOP, and the detection speeds of MobileNetV3 and RegNetY were comparable to YOLOP, this paper only tested the FPS of YOLOP and YOLO-ODL. The experimental results showed that the detection effect and the detection speed of YOLO-ODL were superior to other multi-task models.
Remark 4: The proposed YOLO-ODL model combined the advantages of the improved YOLOv5s model and further improved the accuracy of the model by adopting a weight balance and migration optimization scheme. Therefore, in the aspect of multi-task detection, the proposed YOLO-ODL was superior to other multi-task models. As shown in Figures 7-9, we compared the detection effects of YOLO-ODL and YOLOP on three tasks. Object detection results are shown in Figure 7. YOLOP was prone to false detections and missed detections, mistakenly detecting complex house backgrounds as vehicles, missing the detection of vehicles in the distance. Figure 8 shows drivable area detection results. YOLOP had incomplete detection and many false detection areas, whereas the drivable area detected by YOLO-ODL was more accurate and complete. Lane line detection results are shown in Figure 9. The lane lines detected by YOLO-ODL were more accurate and continuous, but the lane lines detected by YOLOP were discontinuous, and there were more noise areas. TensorRT was used to accelerate the multi-task model and further improve the deployment performance of the model. As shown in Table 6, YOLO-ODL, after the accelerated deployment of TensorRT-FP16, was not only smaller in size but also greatly improved the inference speed. In addition, we also provided Python and C++ APIs to run the model. YOLO-ODL's acceleration ratio after TensorRT-FP16 C++ accelerated deployment reached 5.07, and the detection speed reached 477 FPS, which further verified the efficiency of TensorRT.
Conclusions
In view of the limited computing resources of autonomous vehicles, a multi-task model based on hard parameter sharing was built, which consisted of a shared encoder and three task-specific decoders, which greatly improved computing efficiency. In order to balance the various tasks of the multi-task model, we introduced the loss balance strategy to the multi-task model so that the multi-task model could automatically adjust the weight parameters during training, and we adopted the Mosaic transfer optimization scheme to improve the evaluation index of the multi-task model. Then, the multi-task model was trained and tested to prove the ability of the model to detect the targets, drivable areas, and lane lines and verify the effectiveness of the loss balance strategy and the Mosaic migration optimization scheme. Finally, compared with the latest multi-task network model, it was proved that the proposed YOLO-ODL multi-task model could achieve state-of-the-art performance. TensorRT was used to accelerate the multi-task model and further improve the deployment performance of the model. As shown in Table 6, YOLO-ODL, after the accelerated deployment of TensorRT-FP16, was not only smaller in size but also greatly improved the inference speed. In addition, we also provided Python and C++ APIs to run the model. YOLO-ODL's acceleration ratio after TensorRT-FP16 C++ accelerated deployment reached 5.07, and the detection speed reached 477 FPS, which further verified the efficiency of TensorRT.
Conclusions
In view of the limited computing resources of autonomous vehicles, a multi-task model based on hard parameter sharing was built, which consisted of a shared encoder and three task-specific decoders, which greatly improved computing efficiency. In order to balance the various tasks of the multi-task model, we introduced the loss balance strategy to the multi-task model so that the multi-task model could automatically adjust the weight parameters during training, and we adopted the Mosaic transfer optimization scheme to improve the evaluation index of the multi-task model. Then, the multi-task model was trained and tested to prove the ability of the model to detect the targets, drivable areas, and lane lines and verify the effectiveness of the loss balance strategy and the Mosaic migration optimization scheme. Finally, compared with the latest multi-task network model, it was proved that the proposed YOLO-ODL multi-task model could achieve state-of-the-art performance. | 8,744.4 | 2023-07-01T00:00:00.000 | [
"Computer Science"
] |
Salidroside Protects against Cadmium-Induced Hepatotoxicity in Rats via GJIC and MAPK Pathways
It is known that cadmium (Cd) induces cytotoxicity in hepatocytes; however, the underlying mechanism is unclear. Here, we studied the molecular mechanisms of Cd-induced hepatotoxicity in rat liver cells (BRL 3A) and in vivo. We observed that Cd treatment was associated with a time- and concentration-dependent decrease in the cell index (CI) of BRL 3A cells and cellular organelle ultrastructure injury in the rat liver. Meanwhile, Cd treatment resulted in the inhibition of gap junction intercellular communication (GJIC) and activation of mitogen-activated protein kinase (MAPK) pathways. Gap junction blocker 18-β-glycyrrhetinic acid (GA), administered in combination with Cd, exacerbated cytotoxic injury in BRL 3A cells; however, GA had a protective effect on healthy cells co-cultured with Cd-exposed cells in a co-culture system. Cd-induced cytotoxic injury could be attenuated by co-treatment with an extracellular signal-regulated kinase (ERK) inhibitor (U0126) and a p38 inhibitor (SB202190) but was not affected by co-treatment with a c-Jun N-terminal kinase (JNK) inhibitor (SP600125). These results indicate that ERK and p38 play critical roles in Cd-induced hepatotoxicity and mediate the function of gap junctions. Moreover, MAPKs induce changes in GJIC by controlling connexin gene expression, while GJIC has little effect on the Cd-induced activation of MAPK pathways. Collectively, our study has identified a possible mechanistic pathway of Cd-induced hepatotoxicity in vitro and in vivo, and identified the participation of GJIC and MAPK-mediated pathways in Cd-induced hepatotoxicity. Furthermore, we have shown that salidroside may be a functional chemopreventative agent that ameliorates the negative effects of Cd via GJIC and MAPK pathways.
Introduction
Cadmium (Cd) is a serious environmental toxicant with harmful effects on health in both animals and humans. It is known to target multiple organ systems, particularly the kidneys and liver [1]. Damage resulting from Cd-induced oxidative stress activates signaling cascades, including the Ca 2+ pathway, the mitogen-activated protein kinase (MAPK) pathway, the phosphatidylinositol-3-kinase (PI3K)-Akt pathway and the nuclear factor-κB (NF-κB) pathway, which cause cellular injury, apoptosis and carcinogenesis [2]. However, the definitive signaling pathway that plays the crucial role in Cd-induced apoptosis remains unclear.
Gap junction intercellular communication (GJIC) is one of the most important cellular communications and plays an important role in many biological processes [3,4]. Gap junctions are formed from two connexons on adjacent cells, with each connexon comprised of six connexins (Cx). GJIC, by nature, implies the passive diffusion of small (<1000 Da), hydrophilic substances (e.g., ions, small metabolites and secondary messengers) and there are more than 20 connexin species known to be present in animals and humans, where Cx32 makes up about 90% of the hepatic connexin content. Connexin gene expression and gap junction channel gating are two major mechanisms of GJIC control [5]. It is well known that the functional loss of gap junctions can result in apoptosis, necrosis and carcinogenesis [6][7][8][9] and it is also known that Cd disrupts gap junction activity in hepatocytes in vitro and in vivo [10,11].
MAPKs are a family of Ser/Thr protein kinases of highly conserved enzymes that are unique to eukaryotes. They have been shown to participate in many facets of cellular regulation, such as the control of gene expression, cell proliferation and programmed cell death [12]. Extracellular signal-regulated kinase (ERK), JNK and p38 MAPK are three major members of the MAPK family. These enzymes are activated by phosphorylation and the strength and duration of the activated MAPKs, as well as the cell type, determine the diversity of gene expression, which results in different physiological consequences. Studies have shown that MAPKs are involved in Cd-induced apoptosis in various cell types, including hepatocytes [13].
Salidroside (Sal) is an active constituent of Rhodiola rosea L., which has been used over many years as a medicinal herb for the treatment of altitude sickness [14]. Previous studies have shown that Sal exhibits several pharmacological activities, including anti-oxidant, antiaging, anti-inflammatory, anti-cancer, anti-fatigue and anti-depressant effects [15][16][17][18][19]. Additional studies have found that Sal exerts a protective effect against cellular injury and apoptosis by altering signal transduction in cells. For example, Sal has been shown to protect brain neurons from ischemic injury via the mammalian target of rapamycin (mTOR) signaling pathway [20] and has also been shown to protect against oxygen-glucose deprivation (OGD)/re-oxygenation-induced H9c2 cell necrosis via activation of Akt-Nrf2 signaling [21]. Furthermore, Sal has been found to attenuate H 2 O 2 -induced bone marrow-endothelial progenitor cell (BM-EPC) apoptosis by inhibiting the up-regulation of phosphorylated c-Jun N-terminal kinase (JNK) and p38 MAPK, whilst down-regulating the Bax/Bcl-xL expression ratio [22].
Based on these considerations, in this study we chose Sal as a protective agent against Cd-induced apoptosis to investigate the interactional effects of the MAPK pathway and GJIC, as well as the protective mechanism of Sal both in vitro and in vivo.
Oligonucleotide primers were synthesized by Invitrogen (Shanghai, China). All other reagents were of analytical grade.
Animals and treatment
Thirty female Sprague-Dawley rats weighing 80-100 g were obtained from the Laboratory Animal Center of Jiangsu University (Zhenjiang, China). The animals were housed individually on a 12 h light/dark cycle with unlimited standard rat food and double distilled water (DDW). All experimental procedures were conducted in accordance with the recommendations in the Guide for the Care and Use of Laboratory Animals of the National Research Council and were approved by the Animal Care and Use Committee of Yangzhou University (Approval ID: SYXK (Su) 2007-0005). All surgeries operations were performed under sodium pentobarbital anesthesia, and all efforts were made to minimize any suffering experienced by the animals used in this study. The animals were divided randomly into three groups as follows. (1) Control group: 10 rats consuming DDW as their drinking water. (2) Cadmium group: 10 rats consuming a solution of Cd (50 mg/L) as their drinking water. (3) Cadmium + Sal group: 10 rats treated daily with Sal (35 mg/kg body weight, intragastric gavage, i.g.) and consuming a solution of Cd (50 mg/L) as their drinking water. All rats were sacrificed by cervical dislocation 12 weeks after initial treatment.
Cell culture
Rat liver cells (BRL 3A) were purchased from the Cell Bank of the Institute of Biochemistry and Cell Biology (Shanghai, China). The cells were suspended in DMEM supplemented with 10% FBS, 100 U/mL penicillin and 100 μg/mL streptomycin, and then incubated in a humidified 5% CO 2 /95% air atmosphere at 37°C. BRL 3A cells from passages 10 to 30 were used for all experiments and treated at 90% confluence with different concentrations of Cd (2.5 μM or 10 μM) in the presence or absence of Sal (50 μM) for 12 h.
Real time analysis of cytotoxicity using an xCELLigence DP system The xCELLigence system consists of three components: an analyzer, a real-time cell analyzer (RCTA) station and an E-plate. The dimensionless parameter cell index (CI) was used to quantify cell status and was derived from the measured cell-electrode impedance that directly relates to cell number, cell viability, adhesion and morphology [23].
The xCELLigence system (Roche Applied Science, Basel, Switzerland) was operated according to the manufacturer's instructions [24]. The background impedance of the E-plate was determined in 100 mL medium. Then, 100 mL of the BRL 3A cell suspension was added (10,000 cells/well). Cells were incubated for 30 min in the incubator and then the E-plate was placed into the xCELLigence station. The CI was measured every 15 min. During the phase of rapid CI increase, the growth medium was replaced by serum-free culture medium containing different compounds according to the experimental design. The results were normalized at the end of the assay.
Western blot analysis
Equal amounts of protein (40 μg), were separated by 10% sodium dodecyl sulfate-polyacrylamide gel electrophoresis and transferred onto nitrocellulose membranes. The membrane was incubated with 5% nonfat milk in Tris-buffered saline with 0.1% Tween-20 (TBST) at room temperature for 30 min before incubation with primary antibodies against ERK, P-ERK, JNK, P-JNK, p38, P-p38 (1:1000) or β-actin (1:5000) overnight at 4°C, followed by incubation with horseradish peroxidase (HRP)-conjugated goat anti-rabbit IgG (1:5000) at room temperature for 2 h. The membranes were then washed with TBST and the protein bands were visualized by ECL reagents. The results were analyzed using Image Lab software (Bio-Rad, Hercules, CA, USA).
Scrape-loading/dye-transfer assay for GJIC
A scrape-loading/dye-transfer method (SL/DT) for assessing GJIC using LY was performed as originally described by El-Fouly et al. [25]. In brief, after washing three times with PBS, the cells or fresh liver tissue were scraped with a sharp blade in the presence of LY (0.5 mg/mL), followed by incubation in the dark for 3 min at 37°C and then washed three times with PBS and fixed with 4% paraformaldehyde. The liver tissue samples were then processed by a standard histological technique and mounted on glass slides. The level of GJIC was quantified as the average distance traveled by the LY dye from the scraped edge to the neighboring cells, which was measured using a fluorescent microscope (Leica DMI 3000B, Solms, Germany).
Transmission electron microscopy and immuno-electron microscopy
Internal cellular structures were examined using a Philips CM-100 transmission electronic microscope. Fresh liver tissue was cut into small pieces and fixed with glutaraldehyde (2.5% in 0.1 mol/L cacodylate buffer, pH 7.4) at 4°C for 24 h and embedded in agar. Then, the samples were treated with 1% osmium tetroxide for 2 h. After dehydration with a graded acetone series, the samples were placed in pure acetone and Epon 812 epoxy solution for 30 min and then embedded in Epon 812. Ultrathin sections were cut with a diamond knife and stained with uranyl acetate and lead citrate for transmission electron microscopy. For immuno-electron microscopy, before staining with uranyl acetate and lead citrate, the sections were incubated with Cx32 antibody (1:100) at 4°C overnight, then incubated with 5 nm colloidal gold-conjugated antirabbit IgG (Sigma-Aldrich, Shanghai, China; dilution: 1:100) at 4°C for 4 h.
RNA extraction, reverse transcription and quantitative reverse transcription polymerase chain reaction (qRT-PCR)
Total RNA was extracted from cultured cells and liver tissue using TRIzol Reagent (Invitrogen) according to the manufacturer's instructions. The cDNA was synthesized from 900 ng total RNA using a PrimeScript RT Reagent kit with gDNA Eraser (Takara, Japan). The primers were designed using Primer Premier 5 as follows: β-actin forward, 5'-CGTTGACATCCGTAAAGA CCTC-3' and reverse, 5'-TAGGAGCCAGGGCAGTAATCT -3'; Cx32 forward, 5'-TGGAAG AGGTAAAGAGGCACAAG-3' and reverse, 5'-GGCGGGACACGAAGCAGT-3'. The expression levels of all genes were measured using a real-time PCR system (Applied Biosystems 7500, USA) and the reactions were performed with a SYBR Premix Taq II kit (Takara, Japan) according the manufacturer's instructions. mRNA levels were analyzed with the ΔΔC T method.
Transwell culture of BRL 3A cells BRL 3A cells were seeded at a density of 5×10 4 cells/cm 2 on the inside surface of the polyester membrane of Transwell cell culture inserts (pore size 0.4 μm, surface area 4.67 cm 2 ; Corning, UK), then incubated in serum-free DMEM with 10 μM Cd for 12 h after confluence was achieved. After that, the insert was overturned and BRL 3A cells seeded at a density of 5×10 4 cells/cm 2 on the outside surface of the insert before incubation of the insert in culture medium with 50 μM Sal, 5 μM GA, 10 μM U0126, 10 μM SP600125 or 10 μM SB202190 for 36 h (Fig 1). Finally, cells were carefully wiped from the inside surface of the insert with a cotton swab and the insert washed with cold PBS three times. Cells were fixed in 4% (w/v) paraformaldehyde (Sigma-Aldrich, Shanghai, China) in PBS for 15 min and then stained for 15 min with Hoechst 33258. Images were obtained by an inverted fluorescence microscope (Leica, Germany).
Statistical analysis
All of the experimental data are presented as the mean ± standard deviation (S.D.). Statistical data comparisons among groups were performed using a non-parametric, one way analysis of variance (ANOVA) with p < 0.05 considered statistically significant.
Cd triggers cytotoxicity in BRL 3A cells and rat liver tissue
We evaluated Cd toxicity in BRL 3A cells using the xCELLigence DP system. We found that treatment with Cd (1, 2.5, 5, 10, 20 and 40 μM) resulted in a time-and concentration-dependent decrease of the cell index (CI) in BRL 3A cells (Fig 2A). To assess the influence of Sal on Cd-induced cytotoxicity, cells incubated in the presence of Sal were compared to cells treated with Cd alone. In comparison to BRL 3A cells exposed to Cd alone, the addition of Sal was associated with a slower decrease of the CI (Fig 2B). These results show that Cd triggers cytotoxic injury in a time and dose-dependent manner and Sal at a concentration of 50 μM, can ameliorate the effect of Cd. Transmission electron microscopy was used to confirm that the mitochondria and the nuclei were indeed affected by Cd (Fig 2C). In untreated cells the nuclei were normal, with evenly distributed chromatin, and the mitochondria did not exhibit edema and had clear mitochondrial cristae. Cd was noted to induce nuclei shrinkage and chromatin karyopyknosis, as evidence by the intense staining and marginalization. Mitochondrial swelling was noted and mitochondrial cristae fused partly and became blurry or were even missing. These cellular organelle ultrastructure injuries were partly rescued by treatment with Sal.
GJIC has a dual effect on Cd-induced cytotoxic injury
To determine gap junction activity after treatment of BRL 3A cells and rat liver tissue with Cd, the GJIC was measured by a scrape-loading/dye-transfer method. We found that treatment with Cd resulted in inhibition of GJIC in both BRL 3A cells and rat liver tissue, while Sal had a protective effect on GJIC (Fig 3A). We then investigated the effect of Cd treatment on the mRNA expression level of Cx32 using qRT-PCR. As shown in Fig 3B and Fig 3C, the relative expression level of Cx32 was significantly (p < 0.05 or p < 0.01) decreased following Cd treatment compared with the control group, while co-treatment with Sal significantly attenuated the Cd-induced decrease of Cx32 mRNA. Immuno-labeled Cx32 was observed by transmission electron microscopy ( Fig 3D). In the control group, Cx32 was seen at the gap junctions between adjacent hepatocytes; however, the gold particles appeared decreased in number and scattered after Cd treatment. Notably, Sal appeared to suppress the Cd-induced Cx32 mRNA decrease.
Using a Transwell culture system we co-cultured Cd-exposed cells on the inside surface with BRL 3A cells on the outside surface in the presence of 50 μM Sal and 5 μM GA (Fig 3E). Cd-exposed cells exhibited decreased viability, nuclei chromatin condensation and even nuclear fragmentation. However, these injuries could be ameliorated by the addition of 50 μM Sal. To verify the role of GJIC in Cd cytotoxicity, BRL 3A cells were co-treated with 10 μM Cd and 5 μM GA, which is a prototypical gap junction blocker. As shown in Fig 3F, Cd decreased cell viability and co-treatment with GA exacerbated the reduction in cell viability, as well as nuclei injury, compared with the cells treated only with Cd. However, in the Transwell co-culture system, GA had a protective effect on BRL 3A cells co-cultured with Cd-exposed cells (Fig 3E).
Cd activates the MAPK pathway in BRL 3A cells and rat liver tissue
Western blot analysis was performed to determine the phosphorylation of MAPKs in BRL 3A cells and rat liver tissue (Fig 4A and 4B). The phosphorylation levels of ERK, JNK and p38 in BRL 3A cells were increased after Cd treatment for 12 h, significantly so for ERK. However, Sal co-treatment significantly inhibited the up-regulation of phosphorylated ERK, JNK and p38. Similar results were also seen in rat liver tissue. To define the roles of MAPKs in Cd-induced cytotoxicity in hepatocytes, cells were co-treated with 10 μM Cd and U0126 (10 μM), SP600125 (10 μM) or SB202190 (10 μM). The ERK inhibito (U0126) and the p38 inhibitor (SB202190) prevented the CI decrease induced by Cd (Fig 4C and 4E). However, the JNK inhibitor (SP600125) had little effect on the decrease of the CI induced by Cd (Fig 4B). These data demonstrate that the phosphorylation of ERK and p38 is essential for Cd-induced cytotoxicity and that co-treatment with Sal has a protective effect.
The interaction between GJIC and the MAPK pathway plays an important role in Cd-induced cytotoxicity
To assess the interaction between GJIC and the MAPK pathway following Cd-induced cytotoxicity, MAPK inhibitors (U0126, SP600125 and SB202190) were added to cells in vitro. We found that co-treatment with U0126 significantly blocked decreases in Cd-induced Cx32 mRNA and SB202190 also restored Cx32 levels to some extent, whereas SP600125 had no effect on the expression of Cx32 mRNA (Fig 5A). In addition, alterations in Cx32 mRNA expression were found to be in accordance with GJIC function (Fig 5B). To verify the role of the MAPK pathway in Cd-induced cytotoxicity, BRL 3A cells were incubated in the presence of Cd (10 μM) on the inside surface of a Transwell insert for 12 h (an untreated control was also established). After this time, cells were seeded on the outside surface of the insert in the presence of U0126 (10 μM), SP600125 (10 μM) or SB202190 (10 μM) and compared with untreated control cells and cells treated only with Cd. As shown in Fig 5C, U0126 and SB202190 had a protective effect on Cd-exposed cells, with a decrease in nuclei chromatin condensation and nuclear fragmentation, while SP600125 had little effect. Meanwhile, Fig 5D demonstrates that GA had no effect on the Cd-induced up-regulation of phosphorylated ERK, JNK and p38.
Discussion
Cd is a persistent environmental contaminant with toxic effects in both humans and animals. There is growing evidence that Cd induces apoptosis, but the underlying mechanism remains unclear. The goal of this study was to examine the toxic mechanism of Cd exposure in rat hepatocytes in vitro and in vivo, with a focus on the MAPK pathways and GJIC. In the present study, Cd was demonstrated to be toxic to rat hepatocytes, resulting in decreased cell viability and inhibition of GJIC function. Furthermore, MAPK pathways were found to have critical functions in Cd-exposed hepatocytes, and GJIC inhibition had the dual effect of protecting healthy cells, while damaging injured cells. The protective agent Sal partly attenuated Cdinduced hepatotoxicity. Sal, as an efficacious anti-oxidant, has been widely researched both in vitro and in vivo. A previous study has demonstrated that Sal (10 μM, 50 μM and 100 μM) inhibits 1-methyl-4-phenylpyridinium (MPP+)-induced apoptosis in PC12 cells in a concentration-dependent manner [26]. In addition, Sal (20 or 50 mg/kg) has been shown to reduce neuronal damage in C57BL/6 mice [27]. Based on these reports, we selected Sal at a concentration of 50 μM for our in vitro studies and 35 mg/kg body weight (i.g.), for our in vivo studies. According to our findings, Cd exerted a cytotoxic effect in BRL 3A cells in a time-and concentration-dependent manner (Fig 2A), with ultrastructural damage of nuclei and mitochondria in vivo (Fig 2C). Similar findings have been reported in PC12 cells and rat liver tissue [28][29][30]. These injuries were attenuated by co-treatment with Sal.
In multicellular organisms, the global interplay between the extra-, intra-and inter-cellular communication controls the maintenance of homeostatic balance [31][32][33]. Direct intercellular communication is mainly mediated by gap junctions [34] and the liver was one of the first organs in which GJIC was studied [35]. In the adult liver, gap junctions occupy about 3% of the hepatocyte membrane surface and Cx32 is the major connexin, comprising as much as 90% of the total connexin content [36]. Previous studies have indicated that GJIC can spatially extend apoptosis through the communication of cell death signals from apoptotic cells to healthy cells [8]. In the present study, Cd induced inhibition of GJIC and down-regulation of Cx32 mRNA expression both in BRL 3A cells and the rat liver (Fig 3A, 3B and 3C). Correspondingly, Cx32 decreased and scattered following Cd treatment (Fig 3D). Meanwhile, GJIC inhibition also caused injured cells to lose the rescue signals (such as glucose, ATP and ascorbic acid) provided by healthy cells, which shows a loss of normal growth regulation by the surrounding cells and growth independence [34]. Cells co-treated with GA, a gap junction blocker, exacerbated the effect of Cd in BRL 3A cells (Fig 3F), which is consistent with previous findings, while GJIC inhibition protected healthy cells by limiting the flux of toxic metabolites, such as nitric oxide and superoxide ions, from conjoint damaged cells [37].
To assess whether healthy cells were affected by apoptosis cells, Guo et al. [38] co-cultured normal PC12 cells and Pb 2+ -exposed PC12 cells, which were transfected with the EF1A-eGFP vector. They found that Pb 2+ -exposed PC12 cells induced apoptosis in the unexposed cells via a reactive oxygen species (ROS)-dependent mitochondrial pathway, which was achieved by GJIC. Accordingly, in the present study, Transwell inserts with a 0.4 μm pore size were used in a co-culture system. The pore size was selected so that cells could not cross through the insert membrane but small-size substances could be shared via hemi-channels and/or establishment of gap junctions (as is known for pore sizes less than 3 μm). As such, we co-cultured Cd-exposed BRL 3A cells and normal BRL 3A cells independently on the two sides of a Transwell insert (Fig 1). The resulting findings showed that the inside surface of Cd-exposed cells induced cell damage to the unexposed cells on the outside surface, as shown by nuclei injuries, and that GA and Sal have a protective effect (Fig 3E). This shows that GJIC inhibition has the dual effect of protecting normal cells and exacerbating damage in Cd-exposed cells.
The liver has been identified as a major target of Cd-mediated toxicity; however, not all aspects of the mechanism have been fully elucidated [39]. Cd is known to induce cytotoxic injury via various cell signal-transduction pathways. Wang et al. [40] demonstrated that Cd induced apoptosis via oxidative stress and calcium overload in rat hepatocytes, while Xu et al. [41] have demonstrated that Cd resulted in the caspase-dependent apoptosis of neuronal cells, which was cells. LY transferred to adjacent cells via open gap junctions. Scale bar = 100 μm. ** p < 0.01 compared with control group; # p < 0.05 and ## p < 0.01 compared with the 10 μM Cd group. c: The influence of Cd-exposed cells on BRL 3A cells in a Transwell culture system. Scale bar = 50 μm. d: Effect of GA on the Cd-induced phosphorylation of MAPKs in BRL 3A cells. MAPKs are important signaling enzymes that play a critical role in controlling gene expression, cell survival and cell death; however, the regulation of cytotoxic injury by MAPKs is complex [12]. In this study, we showed that ERK, JNK and p38 activation may be involved in Cd-induced hepatotoxicity both in vitro and in vivo (Fig 4A and 4B), which is in agreement with our previous study [42]. Others studies have shown that ERK is mainly activated by growth factors and tumor promoters and is necesssary for cell proliferation and differentiation, whereas JNK and p38 are involved in apoptosis by promoting cell death [43,44]. Both survival and death signals can activate ERK. Depending on the cell type and stimulus, ERK has a dual effect: besides its involvement in cell proliferation and differentiation, the duration of the activation of ERK also acts as a negative regulator of cell survival and promotes apoptosis [45][46][47]. Our results showed that ERK, JNK and p38 phosphorylation are altered in Cd-exposed BRL 3A cells and the injured liver, and these changes can be attenuated by Sal. Furthermore, ERK and p38 inhibition blocked the decrease observed in the CI following Cd treatment (Fig 4C and 4E). These results showed that ERK and p38 play a crucial role in Cd-induced hepatotoxicity and they are in line with previous findings. However, it should be noted that Chen et al. [48] have developed a different opinion. They found that activation of p38 MAPK is not involved in Cd-induced cell death in PC12 cells, suggesting that p38 may play a different role in different cell types.
In addition, MAPKs are considered to play important roles in GJIC [47,49]. Previous findings have shown that H 2 O 2 -induced GJIC inhibition is involved in both ERK and p38 MAPK activation [50,51], while Cx32 plays a critical role in biological processes of hepatocyte proliferation and cell death [35]. Inhibition of ERK and p38 MAPK-attenuated oxygen-glucose deprivation has been shown to induce Cx32 up-regulation and hippocampal neuron injury [52]. Our findings show that the MAPK inhibitors recover GJIC that was inhibited by Cd ( Fig 5B) and inhibition of ERK and, to some extent, p38 blocked the Cd-induced down-regulation of Cx32 mRNA expression (Fig 5A). In the co-culture system, we observed that U0126 and SB202190 had a protective effect on BRL 3A cells injured through communication with Cd-exposed cells (Fig 5C). Meanwhile, GA, a GJIC blocker, had little effect on the Cd-induced activation of ERK, JNK and p38 (Fig 5D). The results demonstrate that MAPKs induce changes in GJIC most likely by controlling connexin gene expression [53].
In summary, the present work shows that Cd induces rat hepatotoxicity via inhibition of GJIC and activation of MAPK pathways. Interestingly, inhibition of GJIC has the dual effect of protecting healthy cells, while exacerbating injury in damaged cells. ERK and p38 have been found to play critical roles in Cd-induced hepatotoxicity and mediate the function of gap junctions. Finally, Sal may be a potent chemopreventive agent that can prevent the negative effects of Cd via GJIC and MAPK pathways. | 5,930.4 | 2015-06-12T00:00:00.000 | [
"Biology",
"Medicine"
] |
Antiproliferative and Pro-Apoptotic Effect of Novel Nitro-Substituted Hydroxynaphthanilides on Human Cancer Cell Lines
Ring-substituted hydroxynaphthanilides are considered as cyclic analogues of salicylanilides, compounds possessing a wide range of pharmacological activities, including promising anticancer properties. The aim of this study was to evaluate the potential anticancer effect of novel nitro-substituted hydroxynaphthanilides with a special focus on structure-activity relationships. The antiproliferative effect was assessed by Water Soluble Tetrazolium Salts-1 (WST-1) assay, and cytotoxicity was evaluated via dye exclusion test. Flow cytometry was used for cell cycle analysis and detection of apoptosis using Annexin V-FITC/PI assay. Protein expression was estimated by Western blotting. Our data indicate that the potential to cause the antiproliferative effect increases with the shift of the nitro substituent from the ortho- to the para-position. The most potent compounds, 3-hydroxy-N-(3-nitrophenyl)naphthalene-2-carboxamide (2), and 2-hydroxy-N-(4-nitrophenyl)-naphthalene-1-carboxamide (6) showed antiproliferative activity against THP-1 and MCF-7 cancer cells without affecting the proliferation of 3T3-L1 non-tumour cells. Compounds 2 and 6 induced the accumulation of THP-1 and MCF-7 cells in G1 phase associated with the downregulation of cyclin E1 protein levels, while the levels of cyclin B1 were not affected. Moreover, compound 2 was found to exert the pro-apoptotic effect on the THP-1 cells. These results suggest that hydroxynaphthanilides might represent a potential model structure for the development of novel anticancer agents.
Introduction
Salicylanilide derivatives (N-substituted hydroxybenzamides) are known as multitarget agents that possess a wide spectrum of pharmacological activities. These compounds are largely investigated for their promising antibacterial and antimycobacterial effects [1][2][3][4][5]. Some salicylanilides, such as niclosamide or closantel, belong to the class of broad-spectrum anthelmintic agents [6]. Recently, using high-throughput screening, several studies uncovered an antitumor activity of niclosamide, thereby becoming widely studied as a potential anticancer agent [7]. It was proved to effectively induce growth inhibition in a broad spectrum of tumour cell lines together with a minimal toxicity on non-tumour cells [8,9]. On the molecular level, niclosamide inhibited multiple key oncogenic signalling pathways (e.g., Wnt/β-catenin, mTORC1, and NF-κB) [9][10][11][12]. In general, salicylanilide derivatives are presumed to share the structure similarity with the pharmacophore of 4-arylaminoquinazoline derivatives (e.g., gefitinib and erlotinib) that belong to the class of small molecule inhibitors of the protein kinase epidermal growth factor receptor (EGFR PTK) [13][14][15]. This fact led to the intensive research of salicylanilides anticancer properties, as their structure became an attractive model for the design of potent antitumor agents. Several studies were published, in which a series of newly-prepared salicylanilides showed antiproliferative activity against a spectrum of human cancer cell lines, such as promyelocytic leukaemia cells HL-60, chronic myelogenous leukaemia cells K562, human epithelial carcinoma cells A431, or breast carcinoma cells MCF-7. In addition, some salicylanilides have been recently reported to elicit cell cycle arrest or to induce apoptosis in human cancer cell lines [13,[16][17][18].
Recently, several series of various ring-substituted hydroxynaphthanilides were designed and prepared as ring analogues of salicylanilides. Based on the principle of bioisosterism with quinoline-like compounds, the aromatic ring in the salicylanilide pharmacophore was extended by another to obtain the naphthalene scaffold in the structure [3,5]. Compounds containing a quinoline moiety exhibit various pharmacological effects, including anticancer activity [19], hence the hydroxynaphthanilides may possess promising pharmacological properties due to the connection of these two pharmacophores.
The biological activity of salicylanilide pharmacophore could be modified by introducing appropriate substituents in the structure. In addition to a substitution pattern on the salicylic scaffold, SAR studies were focused also on substituents located on the aromatic ring of the anilide part in the structure. It was proved that the biological effects of salicylanilide derivatives are related to both the nature and the position of substituents. The electron parameters of anilide substituents could modify the conformational equilibrium between the closed-ring and open-ring forms of the structure and thus affect the biological activity of the whole molecule. That activity is usually referred to the presence of electron-withdrawing substituents on the anilide moiety [14,20]. In accordance with these findings, our previous results revealed the same relation between the toxicity of ring-substituted hydroxynaphthanilides to the THP-1 cancer cells and the presence of substituents with electron-withdrawing properties [3][4][5]. The SAR studies also found the presence of an electron-withdrawing nitro group to be one of the essential requirements for the anticancer effect of niclosamide [21]. Based on these findings, the substitution by a nitro moiety was determined to be appropriate for the potent anticancer effect of newly-designed hydroxynaphthanilides.
Therefore, we have selected six newly-designed hydroxynaphthanilides, nitro-substituted in different positions on the anilide ring (Table 1), to evaluate their potential anticancer effects in the context of these structural differences. The aim of this work was to assess their antiproliferative activity in two cancer cell lines, THP-1 and MCF-7. Moreover, we also examined the effect on the growth of non-tumour cells 3T3-L1. In addition, changes in cell cycle distribution were evaluated, as well as their pro-apoptotic effect. the class of small molecule inhibitors of the protein kinase epidermal growth factor receptor (EGFR PTK) [13][14][15]. This fact led to the intensive research of salicylanilides anticancer properties, as their structure became an attractive model for the design of potent antitumor agents. Several studies were published, in which a series of newly-prepared salicylanilides showed antiproliferative activity against a spectrum of human cancer cell lines, such as promyelocytic leukaemia cells HL-60, chronic myelogenous leukaemia cells K562, human epithelial carcinoma cells A431, or breast carcinoma cells MCF-7. In addition, some salicylanilides have been recently reported to elicit cell cycle arrest or to induce apoptosis in human cancer cell lines [13,[16][17][18].
Recently, several series of various ring-substituted hydroxynaphthanilides were designed and prepared as ring analogues of salicylanilides. Based on the principle of bioisosterism with quinoline-like compounds, the aromatic ring in the salicylanilide pharmacophore was extended by another to obtain the naphthalene scaffold in the structure [3,5]. Compounds containing a quinoline moiety exhibit various pharmacological effects, including anticancer activity [19], hence the hydroxynaphthanilides may possess promising pharmacological properties due to the connection of these two pharmacophores.
The biological activity of salicylanilide pharmacophore could be modified by introducing appropriate substituents in the structure. In addition to a substitution pattern on the salicylic scaffold, SAR studies were focused also on substituents located on the aromatic ring of the anilide part in the structure. It was proved that the biological effects of salicylanilide derivatives are related to both the nature and the position of substituents. The electron parameters of anilide substituents could modify the conformational equilibrium between the closed-ring and open-ring forms of the structure and thus affect the biological activity of the whole molecule. That activity is usually referred to the presence of electron-withdrawing substituents on the anilide moiety [14,20]. In accordance with these findings, our previous results revealed the same relation between the toxicity of ring-substituted hydroxynaphthanilides to the THP-1 cancer cells and the presence of substituents with electron-withdrawing properties [3][4][5]. The SAR studies also found the presence of an electron-withdrawing nitro group to be one of the essential requirements for the anticancer effect of niclosamide [21]. Based on these findings, the substitution by a nitro moiety was determined to be appropriate for the potent anticancer effect of newly-designed hydroxynaphthanilides.
Therefore, we have selected six newly-designed hydroxynaphthanilides, nitro-substituted in different positions on the anilide ring (Table 1), to evaluate their potential anticancer effects in the context of these structural differences. The aim of this work was to assess their antiproliferative activity in two cancer cell lines, THP-1 and MCF-7. Moreover, we also examined the effect on the growth of non-tumour cells 3T3-L1. In addition, changes in cell cycle distribution were evaluated, as well as their pro-apoptotic effect. the class of small molecule inhibitors of the protein kinase epidermal growth factor receptor (EGFR PTK) [13][14][15]. This fact led to the intensive research of salicylanilides anticancer properties, as their structure became an attractive model for the design of potent antitumor agents. Several studies were published, in which a series of newly-prepared salicylanilides showed antiproliferative activity against a spectrum of human cancer cell lines, such as promyelocytic leukaemia cells HL-60, chronic myelogenous leukaemia cells K562, human epithelial carcinoma cells A431, or breast carcinoma cells MCF-7. In addition, some salicylanilides have been recently reported to elicit cell cycle arrest or to induce apoptosis in human cancer cell lines [13,[16][17][18].
Recently, several series of various ring-substituted hydroxynaphthanilides were designed and prepared as ring analogues of salicylanilides. Based on the principle of bioisosterism with quinoline-like compounds, the aromatic ring in the salicylanilide pharmacophore was extended by another to obtain the naphthalene scaffold in the structure [3,5]. Compounds containing a quinoline moiety exhibit various pharmacological effects, including anticancer activity [19], hence the hydroxynaphthanilides may possess promising pharmacological properties due to the connection of these two pharmacophores.
The biological activity of salicylanilide pharmacophore could be modified by introducing appropriate substituents in the structure. In addition to a substitution pattern on the salicylic scaffold, SAR studies were focused also on substituents located on the aromatic ring of the anilide part in the structure. It was proved that the biological effects of salicylanilide derivatives are related to both the nature and the position of substituents. The electron parameters of anilide substituents could modify the conformational equilibrium between the closed-ring and open-ring forms of the structure and thus affect the biological activity of the whole molecule. That activity is usually referred to the presence of electron-withdrawing substituents on the anilide moiety [14,20]. In accordance with these findings, our previous results revealed the same relation between the toxicity of ring-substituted hydroxynaphthanilides to the THP-1 cancer cells and the presence of substituents with electron-withdrawing properties [3][4][5]. The SAR studies also found the presence of an electron-withdrawing nitro group to be one of the essential requirements for the anticancer effect of niclosamide [21]. Based on these findings, the substitution by a nitro moiety was determined to be appropriate for the potent anticancer effect of newly-designed hydroxynaphthanilides.
Therefore, we have selected six newly-designed hydroxynaphthanilides, nitro-substituted in different positions on the anilide ring (Table 1), to evaluate their potential anticancer effects in the context of these structural differences. The aim of this work was to assess their antiproliferative activity in two cancer cell lines, THP-1 and MCF-7. Moreover, we also examined the effect on the growth of non-tumour cells 3T3-L1. In addition, changes in cell cycle distribution were evaluated, as well as their pro-apoptotic effect.
Effect on Cell Proliferation and Viability
Initially, we examined the effect of six nitro-substituted hydroxynaphthanilides on the proliferation of human leukaemia and breast carcinoma cell lines, using Water Soluble Tetrazolium Salts-1 (WST-1) assay. For such analyses, THP-1 and MCF-7 cells were treated with the compounds at concentrations ranging from 0.5 to 20 µM for 24 h. As shown in Figure 1a, compounds 2, 3, and 6 inhibit cell growth in both cell lines in a dose-dependent manner. The inhibitory effect of 2 and 6 was statistically significant (p < 0.001) starting from the concentration of 2.5 and 5 µM in THP-1 and MCF-7 cells, respectively. From the concentration-response curves, the IC 50 values were determined. As summarized in Table 2, the IC 50 values were found to be 3.06 µM in THP-1 and 4.61 µM in MCF-7 cells for compound 2, and 5.80 and 5.23 µM in THP-1 and MCF-7 cells, respectively, for compound 6. The strongest antiproliferative effect was observed in both THP-1 and MCF-7 cell lines after the treatment with compound 3 (IC 50 1.05 and 1.65 µM, respectively). In contrast, neither compound 1 nor 4 (both ortho-substituted derivatives) was able to induce the inhibition of cell growth in THP-1 or MCF-7 cells at concentrations used in the assay. Compound 5 demonstrated antiproliferative activity only in MCF-7 cells, significant (p < 0.001) at concentrations of 10 and 20 µM (data not shown), however, a 50% reduction in cell growth was not achieved. The proliferation of THP-1 cells was not affected by this compound. After we found that compounds 2, 3, and 6 effectively inhibit the growth of both THP-1 and MCF-7 cancer cells at micromolar concentrations, we assessed additionally their effect on proliferation of non-tumour cell line, 3T3-L1, using WST-1 assay. While compounds 2 and 6 did not decrease cell growth at any of concentrations used, compound 3 affected the proliferation of 3T3-L1 cells in a dose-dependent manner (IC 50 4.41 µM) ( Figure 1b and Table 2).
Subsequently, for the comparison of the antiproliferative and cytotoxic effects we assessed the cell viability after 24 h treatment with compounds 1-6 in both tumour cell lines using the dye exclusion test. In THP-1 cells, we obtained lower LC 50 values: 7.91, 3.44, and 9.98 µM for compounds 2, 3, and 6, respectively. In general, less sensitivity towards the cytotoxic effect of tested compounds was observed in MCF-7 cells. Neither compound 2 nor 6 reduced cell viability under 50% in comparison with the control, while the strongest effect was induced by compound 3 (LC 50 12.91 µM). The results are shown as the means ± standard deviation (SD) of three independent experiments, each performed in triplicate. ** p < 0.01, *** p < 0.001, statistically significant difference in comparison with drug-free control (CTRL).
Effect on Distribution of Cells in Cell Cycle Phases
The cell proliferation assays showed us the ability of selected compounds 2 and 6 to inhibit cancer cell growth. In order to determine at which stage of the cell cycle these compounds induce cell growth inhibition, flow cytometric analyses of cell cycle profiles in THP-1 and MCF-7 cell lines were performed. Cells were exposed to compounds 2 and 6 for 24 h at concentrations exerting significant inhibition of cell proliferation with no or very little concurrent effect on the cell viability. Therefore, THP-1 and MCF-7 cells were treated for 24 h with the compounds at concentrations of 2.5, 5, and 10 μM, respectively. In general, we detected a qualitatively similar effect on the distribution of cells in cell cycle phases following the treatment with compounds 2 and 6 in both leukaemia and breast carcinoma cells. Compounds 2 and 6 induced accumulation of cells in G1 phase in both THP-1 ( Figure 2) and MCF-7 ( Figure 3) cell lines. This was in concert with a simultaneous decrease in the number of cells observed in the S phase compared to the drug-free control, while the percentage of cells in the G2/M phase remained unchanged. The results are shown as the means˘standard deviation (SD) of three independent experiments, each performed in triplicate. ** p < 0.01, *** p < 0.001, statistically significant difference in comparison with drug-free control (CTRL).
Effect on Distribution of Cells in Cell Cycle Phases
The cell proliferation assays showed us the ability of selected compounds 2 and 6 to inhibit cancer cell growth. In order to determine at which stage of the cell cycle these compounds induce cell growth inhibition, flow cytometric analyses of cell cycle profiles in THP-1 and MCF-7 cell lines were performed. Cells were exposed to compounds 2 and 6 for 24 h at concentrations exerting significant inhibition of cell proliferation with no or very little concurrent effect on the cell viability. Therefore, THP-1 and MCF-7 cells were treated for 24 h with the compounds at concentrations of 2.5, 5, and 10 µM, respectively. In general, we detected a qualitatively similar effect on the distribution of cells in cell cycle phases following the treatment with compounds 2 and 6 in both leukaemia and breast carcinoma cells. Compounds 2 and 6 induced accumulation of cells in G1 phase in both THP-1 ( Figure 2 Additionally, the cell cycle analysis allows determining the presence of a subdiploid cell population as a characteristic marker of cells with fractional DNA content. A significant increase (p < 0.001) of the sub-G1 peak was found only after the treatment with 5 µM of compound 2 in THP-1 cells, where an approximately eight-fold increase was observed compared to the drug-free control (Figure 4). In contrast, compound 2 did not induce any elevation of the sub-G1 peak in breast carcinoma cells. Similarly, no significant increase of sub-diploid population of THP-1 or MCF-7 cells caused by 24 h treatment with compound 6 in comparison with the control sample was detected. Next, based on the flow cytometric data that showed the accumulation of cells in the G1 phase upon the treatment with compounds 2 and 6, we examined their effect on the expression of regulatory proteins controlling G1/S and G2/M progression. Whereas total protein levels of cyclin B1 were not changed in THP-1 or MCF-7 cells, the treatment with both compounds 2 and 6 led to the dose-dependent decrease in expression of cyclin E1 (Figures 2c and 3c). Importantly, the levels of cyclin E1 low molecular weight (LMW E1) isoform (42 kDa) were found to be significantly decreased in THP-1 cells. Additionally, the cell cycle analysis allows determining the presence of a subdiploid cell population as a characteristic marker of cells with fractional DNA content. A significant increase (p < 0.001) of the sub-G1 peak was found only after the treatment with 5 μM of compound 2 in THP-1 cells, where an approximately eight-fold increase was observed compared to the drug-free control ( Figure 4). In contrast, compound 2 did not induce any elevation of the sub-G1 peak in breast carcinoma cells. Similarly, no significant increase of sub-diploid population of THP-1 or MCF-7 cells caused by 24 h treatment with compound 6 in comparison with the control sample was detected. Next, based on the flow cytometric data that showed the accumulation of cells in the G1 phase upon the treatment with compounds 2 and 6, we examined their effect on the expression of regulatory proteins controlling G1/S and G2/M progression. Whereas total protein levels of cyclin B1 were not changed in THP-1 or MCF-7 cells, the treatment with both compounds 2 and 6 led to the dose-dependent decrease in expression of cyclin E1 (Figures 2c and 3c). Importantly, the levels of cyclin E1 low molecular weight (LMW E1) isoform (42 kDa) were found to be significantly decreased in THP-1 cells. The results are expressed as the means ± SD of three independent experiments. *** p < 0.001, statistically significant difference in comparison with the drug-free control (CTRL).
Detection of Apoptosis by Annexin V-FITC/PI Assay
To further examine possible pro-apoptotic effect of compound 2 on THP-1 cells, Annexin V-FITC/PI assay was performed using flow cytometry for the quantification of the early and late stages of apoptosis. Staining of cells by Annexin V-FITC conjugate reflects the externalization of phosphatidylserine on the outer surface of the cell membrane as one of the early indicators of apoptosis [22]. In order to obtain further insight into the mechanism of cell death induced by compound 2, we exposed THP-1 cells to a wider concentration range of 2.5, 5, and 10 μM and subsequently analysed the effect at three time-points of incubation (12, 18, and 24 h).
The assay revealed that compound 2 induced a dose-dependent increase of the percentage of early apoptotic as well as late apoptotic/secondary necrotic leukaemia THP-1 cells. In correspondence with the previous detection of a subdiploid cell population compound 2, at concentrations of 2.5 μM, 5 μM, and 10 μM, elicited elevations of Annexin V/FITC-stained cell populations. As shown in Figure 5, this effect was observed even after 12 h of incubation; 10 μM of compound 2 increased significantly (p < 0.01) the proportion of early apoptotic cells to 9.37% in comparison to the percentage of control cells, 2.41%. The same concentration of compound 2 induced the elevation of the number of double-stained cells with incubation time, from 22.48% after 12 h to 41.88% after 24 h incubation. In general, the percentage of late apoptotic/secondary necrotic cells at higher concentrations of compound 2 (5 and 10 μM) prevailed over the early apoptotic cell population at all determined time points.
Detection of Apoptosis by Annexin V-FITC/PI Assay
To further examine possible pro-apoptotic effect of compound 2 on THP-1 cells, Annexin V-FITC/PI assay was performed using flow cytometry for the quantification of the early and late stages of apoptosis. Staining of cells by Annexin V-FITC conjugate reflects the externalization of phosphatidylserine on the outer surface of the cell membrane as one of the early indicators of apoptosis [22]. In order to obtain further insight into the mechanism of cell death induced by compound 2, we exposed THP-1 cells to a wider concentration range of 2.5, 5, and 10 µM and subsequently analysed the effect at three time-points of incubation (12, 18, and 24 h).
The assay revealed that compound 2 induced a dose-dependent increase of the percentage of early apoptotic as well as late apoptotic/secondary necrotic leukaemia THP-1 cells. In correspondence with the previous detection of a subdiploid cell population compound 2, at concentrations of 2.5 µM, 5 µM, and 10 µM, elicited elevations of Annexin V/FITC-stained cell populations. As shown in Figure 5, this effect was observed even after 12 h of incubation; 10 µM of compound 2 increased significantly (p < 0.01) the proportion of early apoptotic cells to 9.37% in comparison to the percentage of control cells, 2.41%. The same concentration of compound 2 induced the elevation of the number of double-stained cells with incubation time, from 22.48% after 12 h to 41.88% after 24 h incubation. In general, the percentage of late apoptotic/secondary necrotic cells at higher concentrations of compound 2 (5 and 10 µM) prevailed over the early apoptotic cell population at all determined time points.
Nevertheless, the different effect was observed after the treatment with two model compounds exerting the pro-apoptotic effect in THP-1 cells. As summarized in Figure 5, cisplatin was found to most effectively increase the rate of early apoptotic cells in a time-dependent manner up to 44.38% after 24 h exposure. While camptothecin increased significantly (p < 0.001) the percentage of both early and late apoptotic cells up to 21.28% and 24.10%, respectively, after 12 h, 24 h treatment led to a decrease of the early apoptotic population to 9.65%; in contrast, late apoptosis increased to 33.08%. Nevertheless, the different effect was observed after the treatment with two model compounds exerting the pro-apoptotic effect in THP-1 cells. As summarized in Figure 5, cisplatin was found to most effectively increase the rate of early apoptotic cells in a time-dependent manner up to 44.38% after 24 h exposure. While camptothecin increased significantly (p < 0.001) the percentage of both early and late apoptotic cells up to 21.28% and 24.10%, respectively, after 12 h, 24 h treatment led to a decrease of the early apoptotic population to 9.65%; in contrast, late apoptosis increased to 33.08%.
Analysis of Proteins Levels Involved in Apoptotic Pathways
Most of the apoptotic signalling pathways are controlled by caspases that belong to a group of cysteine proteases [23]. To assess whether compound 2 affects these signalling cascades and which pathway is activated (intrinsic or extrinsic), the activities of caspase 3, caspase 9, and caspase 8 were evaluated using Western blot analysis. As summarized in Figure 6, after 24 h incubation, compound 2 induced cleavage of pro-caspase 3 dose-dependently; an approximately two-fold decrease of the inactive form upon the treatment with 10 µM compared to the control was detected. Similarly, a comparable two-fold increase of active caspase 3 level was observed after the exposure to the 10 µM concentration of compound 2 in comparison to the control. Additionally, a significant increase of cleaved caspase 9 levels was detected with the most pronounced effects at 10 µM. On the contrary, the level of active caspase 8 was not altered after the treatment with compound 2 in comparison to the control.
Analysis of Proteins Levels Involved in Apoptotic Pathways
Most of the apoptotic signalling pathways are controlled by caspases that belong to a group of cysteine proteases [23]. To assess whether compound 2 affects these signalling cascades and which pathway is activated (intrinsic or extrinsic), the activities of caspase 3, caspase 9, and caspase 8 were evaluated using Western blot analysis. As summarized in Figure 6, after 24 h incubation, compound 2 induced cleavage of pro-caspase 3 dose-dependently; an approximately two-fold decrease of the inactive form upon the treatment with 10 μM compared to the control was detected. Similarly, a comparable two-fold increase of active caspase 3 level was observed after the exposure to the 10 μM concentration of compound 2 in comparison to the control. Additionally, a significant increase of cleaved caspase 9 levels was detected with the most pronounced effects at 10 μM. On the contrary, the level of active caspase 8 was not altered after the treatment with compound 2 in comparison to the control.
Discussion
In the present study, we examined the anticancer effects of a series of newly-synthesized nitro-substituted hydroxynaphthanilide derivatives through the assessment of their antiproliferative activity and cytotoxicity.
Our results showed the difference among the tested compounds in the antiproliferative activity. We found that the potency of cell growth inhibition correlates with the position of the electron-withdrawing nitro group on the anilide ring of the tested compounds. While ortho-substituted derivatives did not elicit any antiproliferative effect in both THP-1 and MCF-7 cancer cells, the shift of the nitro group to the metaor para-position in compounds 2, 3, and 6, led to the cell growth inhibition. Thus, it can be assumed that, most likely, the antiproliferative activity of 3-hydroxynaphthalene-2-carboxanilide and 2-hydroxynaphthalene-1-carboxanilide derivatives increase depending on the position of the nitro group as follows: ortho < meta < para. This different activity could be possibly related to the steric effect of the anilide substituents. Recently, it was described that the presence of a substituent in the ortho position causes the twist of the whole aniline ring plane towards the naphthalene scaffold, while metaand, especially, para-substituted derivatives have a practically linear molecule [24]. Moreover, not only the location of the substituent on the anilide moiety but also the position of the β-ring of naphthalene towards the phenolic and carboxanilide moietis affected the intensity of the antiproliferative effect of these compounds. In our study, stronger antiproliferative activity was observed in substituted 3-hydroxynaphthalene-2-carboxanilides when comparing the IC 50 values of meta-substituted compounds 2 and 4 or para-substituted 3 and 6. The similar structure-activity relationship was determined for the cytotoxicity of the tested compounds. Nevertheless, compounds 2, 3, and 6 exerted stronger antiproliferative rather than cytotoxic effect in cancer cells; approximately 2-3-fold higher LC 50 values compared to IC 50 values were obtained in the assays on THP-1 cells. Even more pronounced difference was observed in MCF-7 cells, where the LC 50 values were achieved only upon the treatment with compound 3, with an approximately seven-fold higher dose in comparison with IC 50 .
To assess whether tested compounds also influence the growth of other than cancer cells, we have extended our antiproliferative analysis and employed non-tumour fibroblast cell line 3T3-L1. Compound 3 that exerted the most substantial antiproliferative and cytotoxic effects towards both cancer cell lines was also capable of inhibiting the growth of the non-tumour line. Interestingly, a different effect was observed upon the treatment with compounds 2 and 6, where such antiproliferative activity in non-tumour cells was not detected. Results of antiproliferative effects showed us that among all tested compounds, compounds 2 and 6 were the most potent and, thus, were chosen for further, more detailed analyses.
One characteristic feature of cancer cells is the deregulation of the cell cycle, which leads to their uncontrolled proliferation. Therefore, the inhibition of cell cycle progression represents a common target of anticancer agents [25]. We performed the cell cycle analysis to reveal whether the antiproliferative effect of compounds 2 and 6 is reflected in the modification of cell cycle progression. Our results showed that both compounds were able to accumulate THP-1 and MCF-7 cancer cells in the G1 phase and to inhibit the transition of cells to the synthetic phase. We assume that this most likely reflects the antiproliferative effect observed in both cell lines (Figure 1a). The progression through the cell cycle is mediated by a family of cyclin-dependent kinases, the activity of which depends on the binding of the regulatory proteins, cyclins [26]. The observed accumulation of THP-1 and MCF-7 cells in the G1 phase after the treatment with compounds 2 and 6 was accompanied by a reduction of cyclin E1 level in a dose-dependent manner ( Figure 2c). As the activator of CDK2, cyclin E1 is responsible for the G1/S phase progression and, thus, it is involved in surpassing the restriction point [27]. Many cancers typically overexpress cyclin E1, which is also proved in the MCF-7 cell line [28]. This might support our finding of only slight downregulation of cyclin E1 caused by the treatment of MCF-7 cells with compounds 2 and 6, although these compounds effectively inhibited the G1/S transition. Interestingly, besides the downregulation of cyclin E1 full-length form, we also detected a more pronounced reduction of LMW E1 isoform levels in THP-1 cells treated with compounds 2 and 6. LMW E1 isoforms are generated primarily in cancer cells, where they still remain fully functional. They have even higher potency to increase CDK2/E1 activity than the full-length form and, thus, they move the cells through the cell cycle more effectively than the full-length form [29,30].
Our previous study reported a similar detection of the decreased levels of cyclin E1 isoforms in THP-1 cells treated with geranylated flavanone tomentodiplacone B that coincided with an induced accumulation of cells in G1 phase [31]. While cyclin B1 is involved in the G2/M transition associated with CDK1 [26], we did not observe any change in the levels of cyclin B1 in THP-1 or MCF-7 cells after the exposure to compounds 2 and 6. These findings are supported by our flow cytometric data that did not indicate any significant difference in the proportion of cells in G2/M cell cycle phase upon the treatment with these compounds (Figures 2b and 3b). Based on those results, we could suggest that compounds 2 and 6 most likely affect G1/S rather than the G2/M transition.
The presence of cell nuclei with hypodiploid DNA content during the cell cycle analysis could indicate a possible presence of apoptotic cells [32]. The assessment of sub-G1 peak levels revealed different effects among the tested compounds; a significant increase was detected only in THP-1 cells upon the treatment with compound 2 (Figure 4). Based on these findings, we performed further analysis to prove its possible pro-apoptotic effect in the THP-1 cell line. Results of Annexin V-FITC/PI assay showed us that compound 2 induced the THP-1 cells to undergo an early stage of apoptosis even after 12 h exposure ( Figure 5). Nevertheless, compound 2 accumulated more effectively (dose and time-dependently) in cells in the late apoptotic stage. These results correlate with the data obtained from the viability staining assay. In addition, two already known anticancer agents of a different mode of action, cisplatin, which is able to crosslink with the DNA and, thus, cause DNA damage [33], and camptothecin as the S-phase-specific inhibitor of the enzyme DNA topoisomerase-I [34], were added to the assay as model compounds with proved pro-apoptotic effects in THP-1 cells [35,36]. Although our results found all three compounds to significantly increase the number of cells positive for Annexin V-FITC staining, their effect led to different proportions of early and late apoptotic/secondary necrotic cells. While cisplatin induced a time-dependent substantial increase in the fraction of early apoptotic cells, camptothecin most likely elicited the time-dependent transfer of cells from early apoptotic to late apoptotic stages. These differences observed in the effect of three tested compounds enable us to presume a different mechanism of action of compound 2 in comparison with one of the two model anticancer agents.
These findings prompted us to further investigate the involvement of compound 2 in the apoptotic pathways. The caspases regulate the process of apoptosis in a different manner. The activation of caspase 8 is realized through the extrinsic apoptotic pathway after the binding of a ligand to an appropriate death receptor. Subsequently, the active form interacts with effector caspase 3 and that results in its cleavage and activation. On the other hand, initiator caspase 9 is involved in the intrinsic, also known as the mitochondrial apoptosis pathway, and is activated after the leakage of the mitochondrial cytochrome c. This also leads to proteolytic cleavage of inactive procaspase 3 and to its activation. Therefore, it denotes the essential role of caspase 3 in both extrinsic and intrinsic pathways, as it also comprises a link between them [37,38]. After 24 h treatment, compound 2 was found to be capable of inducing an increase of active caspase 3 level, including the decreased level of inactive pro-caspase 3, both significantly at a concentration of 10 µM (Figure 6). At the same time, compound 2 caused also the cleavage of pro-caspase 9. On the contrary, no change in the level of the active form of caspase 8 was observed in comparison with the control, non-treated cells. These results indicate that compound 2 induces apoptosis in THP-1 cells by activating a caspase cascade. In addition, we could hypothesize that this compound might be preferably involved in the intrinsic apoptotic pathway. However, such specificity needs to be proved by additional analyses, and the mechanism of targeting apoptotic pathway remains unknown.
Chemicals and Reagents
The tested nitro-substituted hydroxynaphthanilides 1´6 were prepared and supplied by the Department of Chemical Drugs, Faculty of Pharmacy, University of Veterinary and Pharmaceutical Sciences Brno, Czech Republic. The synthesis and structural characterization of these compounds have been described previously [3,5]. Due to poor solubility in water, the compounds were dissolved in dimethyl sulfoxide (DMSO) (Sigma-Aldrich, St. Louis, MO, USA), while the stock solutions were prepared freshly before each experiment. The final concentration of DMSO in the assays never exceeded 0.1% (v/v). Cisplatin and camptothecin were purchased from Sigma-Aldrich. RPMI 1640 and DMEM culture media, phosphate-buffered saline (PBS), foetal bovine serum (FBS) and antibiotics (penicillin and streptomycin) were obtained from HyClone Laboratories, Inc. (GE Healthcare, Logan, UT, USA). Mouse monoclonal antibodies against cyclin E1 (sc-247), caspase 3 (sc-7272) and caspase 9 (sc-17784) were purchased from Santa Cruz Biotechnology (Santa Cruz, CA, USA). Rabbit polyclonal antibodies against cyclin B1 (ab2949) and caspase 8 (ab-25901) were purchased from Abcam (Cambridge, UK). All other reagents, unless specified elsewhere, were purchased from Sigma-Aldrich.
Cell Culture
THP-1 human monocytic leukemia cell line, MCF-7 human breast adenocarcinoma cells and 3T3-L1 mouse embryonic fibroblast were purchased from the European Collection of Cell Cultures (ECACC, Salisbury, UK). Cells were routinely tested for the absence of mycoplasma (Hoechst 33258 staining method). THP-1 cells were maintained in RPMI 1640 culture medium containing 2 mM L-glutamine; MCF-7 and 3T3-L1 cells were cultured in DMEM medium. All of the culture media were supplemented with 10% heat-inactivated FBS and antibiotics (100 U/mL penicillin and 100 µg/mL streptomycin). Cells were maintained at 37˝C in a humidified atmosphere containing 5% CO 2 .
Cell Cycle Analysis
THP-1 and sub-confluent MCF-7 cells were treated and subsequently incubated with indicated concentrations of compounds 2 and 6 for 24 h. After the incubation, cells were washed twice in PBS (pH 7.4), fixed in 70% ethanol and stored at´20˝C overnight. Fixed cells were collected by centrifugation, and supernatant was discarded. The cell pellet was washed twice with PBS and incubated with RNaseA (0.02 mg/mL) and 0.05% (v/v) Triton X-100 in PBS for 30 min at 37˝C. After the nuclei staining with propidium iodide (PI) (0.04 mg/mL), the cell cycle distribution was analysed using a flow cytometer Cell Lab Quanta SC (Beckman Coulter, Brea, CA, USA). The quantification of cell cycle distribution was carried out using software MultiCycle AV (Phoenix Flow System, San Diego, CA, USA). A total number of 2ˆ10 4 cells was analysed per sample.
Detection of Apoptosis Using Annexin V-FITC/PI Assay
Early and late stages of apoptosis were detected using Annexin V-FITC Kit-Apoptosis Detection Kit faccording to the manufacturer´s instructions. THP-1 cells were treated with increasing concentrations of compound 2 (2.5, 5, and 10 µM), cisplatin (10 µg/mL) and camptothecin (5 µM). At each time-point of incubation (12,18, and 24 h) the cells were washed with ice-cold PBS prior to being resuspended at a concentration of 5ˆ10 6 cells/mL in a total volume of 100 µL of 1ˆbinding buffer. Annexin V-FITC solution (final concentration 0.25 µg/mL) and PI (final concentration 12.5 µg/mL) were added to each sample; the cell suspension was kept on ice and incubated for 15 min in the dark. After that, the analysis was carried out by flow cytometry. The data were evaluated using Kaluza Flow Cytometry Analysis 1.2. Per sample, a total number of 2ˆ10 4 cells were analysed.
Western Blotting
For Western blotting, cells were washed with PBS and lysed in lysis buffer (100 mM Tris-HCl, pH = 6.8; 20% glycerol; 1% SDS) containing protease and phosphatase inhibitor cocktails. Protein concentration was measured using Roti ® -Quant universal (Carl Roth, Karsruhe, Germany) according to the manufacturer's instructions. Cell lysates were supplemented with bromophenol blue (final concentration 0.01% (w/v)) and β-mercaptoethanol (final concentration 1% (v/v)) prior to being heated for 5 min at 95˝C. Equal amounts of protein (10 µg) were loaded into a 12% polyacrylamide gel, separated by SDS-polyacrylamide gel electrophoresis and subsequently electrotransferred onto nitrocellulose membranes. Reversible Ponceau S. staining was performed to assess equal sample loading. Then, the membranes were blocked with 5% non-fat dry milk in TBST (10 mM Tris-HCl pH = 7.5, 150 mM NaCl, 0.1% (v/v) Tween-20) and appropriate primary and secondary antibodies were used for immunodetection. The proteins were visualized by ECL Plus reagent according to the manufacturer's instructions. The intensity of bands was semi-quantitatively analysed using the ImageJ software (National Institute of Mental Health, Bethesda, MD, USA).
Statistical Analysis
All experimental data were expressed as the arithmetical mean˘standard deviation (SD). Statistical analysis was performed using one-way analysis of variance (ANOVA) followed by the Dunnett's post test using GraphPad Prism 5.00 software. Statistical significance was assessed at levels of p < 0.05, p < 0.01 and p < 0.001.
Conclusions
The present study provides the first description of the antiproliferative activity of nitro-substituted hydroxynaphthanilides in the context of structure-activity relationships. Our results indicate that the potency of ring-substituted hydroxynaphthanilides towards cell growth inhibition increases with positioning of the nitro group as follows: ortho < meta < para. The most promising compounds 2 and 6 exerted antiproliferative activity in THP-1 and MCF-7 cancer cells with single-digit micromolar IC 50 values, while they had a minimal effect on the growth of 3T3-L1 non-tumour cells. Compounds 2 and 6 accumulated cancer cells THP-1 and MCF-7 in G1 cell cycle phase, which was accompanied by the observed down-regulation of cyclin E1 levels. Moreover, compound 2 was found to induce apoptosis in THP-1 cells via a caspase-mediated cascade. The results also indicate that apoptosis was probably induced through the intrinsic apoptotic pathway, although further analysis is still required to verify such assumption. According to the results, nitro-substituted hydroxynaphthanilides 2 and 6 can be considered as potential anticancer agents, and the structure of hydroxynaphthanilides is an appropriate model moiety for further design of compounds with potential anticancer properties. | 9,067.8 | 2016-07-28T00:00:00.000 | [
"Chemistry",
"Biology"
] |
Mesoscopic spin transport between strongly interacting Fermi gases
We investigate a mesoscopic spin current for strongly interacting Fermi gases through a quantum point contact. Under the situation where spin polarizations in left and right reservoirs are same in magnitude but opposite in sign, we calculate the contribution of quasiparticles to the current by means of the linear response theory and many-body T -matrix approximation. For a small spin-bias regime, the current in the vicinity of the superfluid transition temperature is strongly suppressed due to the formation of pseudogaps. For a large spin-bias regime where the gases become highly polarized, on the other hand, the current is affected by the enhancement of a minority density of states due to Fermi polarons. We also discuss the broadening of a quasiparticle peak associated with an attractive polaron at a large momentum, which is relevant to the enhancement.
I. INTRODUCTION
Quantum simulation with ultracold atomic gases allows one to explore regimes of quantum-many body problems where conventional systems such as condensed matter and nuclear matter are hard to reach [1][2][3]. A strongly-interacting Fermi gas realized with the Feshbach resonance is the prototype example, and revealed existences of Bardeen-Cooper-Schrieffer (BCS) to Bose-Einstein condensation (BEC) crossover and the unitarity regime in which the typical length scale characterizing the atomic interaction disappears [4][5][6]. Whilst both theoretical and experimental progresses have deepened understanding of the bulk phase structure, understanding of the non-equilibrium properties remains challenging.
Recently, quantum transport of strongly-interacting Fermi gases attracts rising attention in association with the atomtronics devices where non-equilibrium properties on circuit or two-terminal systems are investigated [7]. By using the controllable ultracold Fermi gases, quantum point contact [8], and junction systems attached with two [9] or three [10] dimensional reservoirs have experimentally been implemented. In the case of the stronglyinteracting superfluid regime, superfluid transport such as nonlinear current-bias characteristics induced by the multiple Andreev reflections [11], and the AC Josephson oscillation [10] and DC Josephson effect [12] has been confirmed. In the case of the normal Fermi gas with a strong attractive interaction, a conductance beyond the quantized value has been found [13].
In addition to controllability of the interaction, transport systems with ultracold Fermi gases have an advantage that spin transport can directly be measured with a tunable spin bias [13]. We note that this is in contrast to condensed matter systems where spin transport is normally measured in an indirect manner [14]. In the context of the interacting Fermi gases, while mass transport * These two authors contributed equally to this work. contains information both on single-particle and collective excitations, spin transport is not involved in Cooper pair transport [15,16]. Thus, it is expected that the spin transport measurement becomes a sensitive probe of single-particle excitations.
In this paper, we investigate mesoscopic spin transport in attractively interacting two-component Fermi gases. We consider the tunneling regime of a two terminal system consisting of two normal Fermi gases through a quantum point contact (see Fig. 1). Spin polarizations of the Fermi gases in left and right reservoirs are assumed to be equal in magnitude but opposite in sign and are controlled by the parameter h. Due to this spin imbalance, spin-up and spin-down fermions can move in opposite directions, and therefore a nonzero spin current without flow of a mass current is generated. We focus on the strongly interacting regime near unitarity, where the absolute value of a s-wave scattering length |a| is much larger than the interatomic distance, and intercomponent fermions strongly interact with each other. Figure 2 is the phase diagram of the spin-imbalanced Fermi gases at unitarity [17,18]. When the low temperature and small h regime is concerned, each reservoir becomes the superfluid where spin excitations are expected to be frozen. Therefore, in this work, we elucidate spin transport above the superfluid critical temperature.
In contrast, excitations of a highly polarized Fermi gas realized in the large h regime is governed by Fermi polarons, which are mobile impurities immersed in a Fermi sea [49]. Polaronic properties such as renormalization factors, effective masses, and polaron energies have been measured with RF spectroscopy [50][51][52][53][54][55][56][57]. Correspondingly, a bunch of theoretical works have been made [58][59][60][61][62][63][64][65][66][67][68][69][70][71][72][73][74][75][76][77], most of which consider the zero-temperature and single-polaron limit by assuming that the impurity chemical potential is equal to the polaron energy. However, we note that the theory based on such an ideal limit is not available to discuss polaronic effects in spin transport at finite temperature and given chemical potential. To incorporate the polaronic properties in a correct fashion, we perform the finite-temperature many-body formalism of Fermi polarons. By using this formalism, moreover, we show that the crossover from the pseudogap to the polarons can be explored through spin transport.
This paper is organized as follows. In Sec. II, we present the formalism of the tunneling Hamiltonian approach together with the diagrammatic T -matrix approximation. Section III is devoted to discussing how excitation properties in strongly interacting Fermi gases in reservoirs affect spin transport. We conclude this paper in Sec. IV.
Throughout this paper, we set = k B = 1 and the volumes for both reservoirs are taken to be unity.
II. FORMALISM
In order to study spin transport of two terminal systems in normal Fermi gases with strong interaction, we begin with the following grand canonical Hamiltonian: where left and right reservoirs are referred to as j = L, R, respectively, c p,σ,j is the annihilation operator of a fermionic atom with spin σ = ↑, ↓ in the reservoir j, and ξ p,σ,j = p 2 /2m − µ σ,j is a single-particle energy measured from the chemical potential µ σ,j . The interaction between fermions in the reservoir j is attractive (U > 0) and related to the s-wave scattering length a by 1/U = p m/p 2 − m/(4πa). The term in H T denotes the tunneling of fermions from one reservoir to the other and t characterizes the strength of the tunneling. As mentioned above, we focus on the situation where there is a pure spin bias shown in Fig. 1. The majority and minority components in the left reservoir are ↑ and ↓, respectively, and those in the right reservoirs are the other way around. It follows that the parameter h controls not only the polarizations in both reservoirs but also a bias in spin transport. We assume that both reservoir has same temperature T , which is above the superfluid transition temperature. For convenience, we define a new label α ≡ (σ, j) and, hereafter, the majority components (σ, j) = (↑, L), (↓, R) are referred to as "α = +," while the minority ones (σ, j) = (↓, L), (↑, R) as "α = −." The chemical potentials of the majority (α = +) and minority (α = −) are µ ± = µ ± h as depicted in Fig. 1.
The spin current operator in the Heisenberg representation is given byÎ p,σ,j c p,σ,j is the particle number operator. By using the linear response theory, the spin current to leading order in the tunneling amplitude t can be obtained in a similar way as for a mass current [78]. For a steady state, the spin current is obtained as (see Appendix A) where ρ α (ω) is the DOS for the majority (α = +) [minority (α = −)] and f (ω) = 1/(e ω/T + 1) is the Fermi distribution function. We note that regardless of value of h, Eq. (2) is correct up to t 2 . In order to take contributions of pair correlations to ρ α (ω) into account, we employ the extended T -matrix approximation (ETMA) [17,71]. The density of states is related to an analytically-continued Matsubara Green's function G α (p, iω n ) with ω n = (2n + 1)π T (n ∈ Z) as follows: where ξ p,α = p 2 /2m − µ α , and δ is an infinitesimal positive constant. Within the ETMA, the self energy Σ α (p, iω n ) and the many-body T -matrix Γ(q, iν ) with ν = 2 π T ( ∈ Z) are given by Feynman diagrams in Fig. 3, leading We note that the self-consistent programme above can reduce unphysical results. For example, the ordinary Tmatrix approximation is known to suffer from a negative spin susceptibility in the strong-coupling regime. However, ETMA spin susceptibility takes a positive value in the whole crossover regime [17]. We now discuss the choice of parameters (a, T, µ, h) in this work. In an unpolarized case (h = 0), we fix (a, T, µ) as follows. For a given density n 0 of the total particle number in each reservoir, the corresponding Fermi momentum and temperature are provided by k F,0 = (3π 2 n 0 ) 1/3 and T F,0 = k 2 F,0 /(2m), respectively, and two dimensionless parameters (k F,0 a) −1 and T /T F,0 are fixed. Then, (k F,0 a) −1 → −∞ corresponds to the weak interaction limit and (k F,0 a) −1 → +∞ to the strong interaction limit in the fermion language. The spinaveraged chemical potential µ is determined so that the particle number equation in the absence of h, is satisfied. Then, the fictitious Zeeman field h is changed with (a, T, µ) fixed. Since we are interested in how strong correlations affect spin transport, we consider the regime |k F,0 a| −1 1 near unitarity. Figure 2 shows the phase diagram of the Fermi gases at unitarity ((k F,0 a) −1 = 0) in the (h, T )-plane. The transition temperature is determined as the temperature satisfying the Thouless criterion [79] given by [Γ(0, 0)] −1 = 0. We note that at low temperature there is the first-order phase transition between normal and superfluid phases, including the socalled Fulde-Ferrel-Larkin-Ovchinnikov state [80]. To address these transitions, one has to calculate the thermodynamic potential in each phase, which is out of scope in this paper. Figure 4 shows the number densities n ± of majority and minority atoms at unitarity as functions of h. (Note that the total number density n + + n − changes as h increases.) The monotonic behavior of x = n − /n + in Fig. 4 means that, as h becomes larger, the gases in both reservoirs become more highly polarized. As shown in Figs. 10 and 11 in the next section, the polaron picture becomes valid for the large h regime.
III. RESULT
We now discuss spin transport properties for normal Fermi gases with strong interparticle interactions under the configuration shown in Fig. 1. Figure 5 shows the spin conductance G spin ≡ I spin /(2h) at unitarity. We can see that G spin grows with increasing h with a fixed T . Furthermore, it is remarkable that G spin is largely suppressed in the low-T and low-h regime, where the pseudogap emerges. As shown in Fig. 6, our ETMA calculation for unpolarized (h = 0) and polarized (h = 0.25ε F,0 ) gases at T = T c can confirm the pseudogap structures of DOSs, being the signature of the preformed Cooper pairs in the BCS-BEC crossover regime of an ultracold Fermi gas.
To make the effect of the pseudogap on spin transport clearer, let us focus on the zero-bias limit. In this limit, majority and minority DOSs become identical,ρ(ω) = ρ ± (ω)| h=0 , and the spin conductance reduces to the following form: limit. We note that the calculation of G spin is stopped at the transition temperature T c . Away from T c , G spin increases with decreasing T because of quantum degeneracy of fermions. On the other hand, as shown in Fig. 6, the DOS has a dip structure around ω = 0 near the superfluid transition. Since −∂f (ω)/(∂ω) ∝ cosh −2 (ω/2T ), the spin conductance is sensitive toρ(ω) around ω = 0. Therefore, the appearance of this pseudogap leads to a large suppression of G spin . The single-particle excitation is strongly suppressed due to the formation of spin-singlet pairs in the pseudogap regime. This suppression leads to the so-called spin gap in the temperature dependence of the spin susceptibility [20,[34][35][36][37]39]. We note that the spin-gap temperature, where the spin susceptibility starts to be suppressed due to strong pairing fluctuations, is T SG = 0.37T F,0 at unitarity [39]. Although it is the crossover temperature and there are ambiguities for the definition to characterize these phenomena, the maximum temperature of G spin at unitarity is also close to T SG . This result indicates that G spin is also useful to study pseudogap physics. From Fig. 7, we can also see the interaction dependence of G spin . In the weak-coupling side ((k F,0 a) −1 = −0.5 in Fig. 7), G spin becomes larger compared to that at unitarity. However, even at this coupling, G spin near T c decreases due to the pairing fluctuation effects as decreasing temperature. As in the case of the spin susceptibility [39], the spin conductance is strongly suppressed as increasing strength of the pairing interaction. At stronger couplings, the gases in both reservoirs are dominated by tightly bound molecules and the spin degree of freedom tends to be frozen. In this case, only the thermally dissociated atoms contribute to the spin susceptibility as well as spin transport. Figure 8 represents a crossover of the spin current at unitarity from the pseudogap regime to the polaronic regime by shifting h. We can see that I spin for small h is smaller than in the non-interacting counterpart, where the current is analytically given by Eq. (B2) in Appendix B. As explained in the discussion of G spin , this suppression is caused by the pseudogap in the region where h is small. Figure 9 shows the calculated DOSs for both majority and minority components at T = 0.21T F,0 for various h. When h becomes larger, polarizations of the gases in both reservoirs grow and the pseudogap structures of ρ ± (ω) vanish since the gases go away from T c at a fixed temperature. The majority DOS is enhanced in the whole energy region with increasing h due to the increase of n + . In the large-h region, ρ + (ω) coincides with a DOS in an ideal Fermi gas given by Eq. (B1), since under a large population imbalance, minority atoms cause a negligible effect to a large number of majority atoms. On the other hand, the minority DOS shows a more complex modification than ρ + (ω) with increasing h. In particular, in the large-h regime, minority atoms can be regarded as the so-called Fermi polarons. In our configuration of chemical potentials with fixed µ = (µ + + µ − )/2, I spin is enhanced compared with that without self-energy corrections [81]. This implies that the polaronic quasi-particle excitations encoded in the self-energy corrections play an important role in spin transport for a large spin bias.
In order to discuss the contributions of the polaronic transport to I spin for large h, we start with the investigation of the single-particle spectral functions defined by In the literature on the Fermi polarons, the Fermi momentum k F,+ = (6π 2 n + ) 1/3 and the Fermi energy ε F,+ = k 2 F,+ /(2m) for majority atoms are conventionally taken as units of energy and momentum. Thus we use these units to discuss A α (p, ω) in the large spin-bias regime. For large h, the majority spectral function can be replaced by that for free fermions A + (p, ω) = δ(ω +µ + − p 2 2m ). On the other hand, Fig. 10 shows that the minority spectral functions near unitarity are distinct from its non-interacting counterpart. In particular, there are two types of characteristic excitations: a sharp peak at low energy corresponding to an attractive polaron and a broad peak around ω + µ − ≈ ε F,+ associated with a repulsive branch or repulsive polarons. Since A − (p, ω) at low energy is dominated by the attractive polaron, it is well approximated by where Z a , m * a , and E a are the renormalization factor, the effective mass, and the energy of the attractive polaron, respectively. By integrating this approximated form over p, the contribution from the polaron to ρ − (ω) is found to be for ω > −µ − + E a . Figure 11 shows the minority DOS at unitarity numerically calculated in the ETMA. The obtained ρ − (ω) for ω + µ − ε F,+ is found to be enhanced compared to a DOS without self-energy corrections. Fitting the calculated DOS at low energy by using Eq. (11), we obtain Z a (m * a /m) 3/2 = 0.945 and E a = −0.627ε F,+ . These results are in good agreement with the zerotemperature results in the single-polaron limit, where Z a = 0.78, m * a /m = 1.17, and E a = −0.606ε F,+ , leading to Z a (m * a /m) 3/2 = 0.987 [59,64]. This means that the effects of finite T and x on ρ − (ω) are not so important in this temperature and bias range [71]. For 0 ω + µ − ε F,+ , ρ − (ω) deviates from Eq. (11). In this energy range, the enhancement of ρ − (ω) comes from not only the attractive polaron at high momentum, whose lifetime is finite, but also the repulsive branch. In particular, the broadening of the peak associated with the attractive polaron in Fig. 10 can intuitively be understood as follows. The polaron is a quasiparticle consisting of a minority atom surrounded by majority atoms. This quasiparticle picture is valid when the velocity of a dressed minority atom v − ≈ p/m * a is smaller than a typical velocity of the majority atoms v F,+ = k F,+ /m. When v − > ∼ v F,+ , the majority atoms can no longer follow the fast-moving minority atom and thus the attractive polaron tends to be unstable. At unitarity, this unstable regime is estimated as p > ∼ (m * a /m)k F,+ ≈ 1.17k F,+ and is consistent with the region where the peak associated with the attractive polaron becomes broad (see Fig. 10). This mechanism of the broadening is analogous to the Cherenkov instability in Bose polarons [82,83], where a minority atom undergoes the supersonic regime associated with the Bogoliubov phonons. Actually, this intermediate energy range of polaron spectra plays a significant role to understand polaronic spin transport as discussed below. For sufficiently large ω, ρ − (ω) becomes close to its non-interacting counterpart, which is consistent with the asymptotic behavior derived by the operator product expansion [84].
Let us now come back to I spin for a large spin bias. Figure 12(a) shows functions in the integrand of Eq. (2) at T = 0.25T F,0 and h = ε F,0 . We can see that ρ + (ω − h) is almost consistent with that without self-energy corrections. Because of the existences of f (ω − h) − f (ω + h) reflecting the Fermi-Dirac statistics and of a threshold energy ω = −µ for ρ + (ω − h), ρ − (ω + h) contributes to I spin only over the region where the attractive polaron at finite momentum and the repulsive branch appear. We note that in the absence of the pairing interaction, both curves coincide with each other. Thus, the enhancement of the spin current in the highly polarized regime originates from the polaron excitations in the intermediateenergy range (0 ω + µ − ε F,+ in Fig. 11). Figure 12(b) compares I spin with and without selfenergy corrections at (k F,0 a) −1 = −0.5, 0, and 0.5. While the enhancement of I spin is large in the weak-coupling side, it is small in the strong-coupling side. Since the magnitude of I spin depends on µ obtained by solving Eq. (7) at each (k F,0 a) −1 in our configuration, I spin becomes larger in the weak-coupling side where µ is positively large compared to that in the strong-coupling side. In this sense, the polaron properties appear as the ratio given by I spin with and without the self-energy corrections. This ratio is still larger in the weak-coupling side. The enhancement of I spin is associated with the overlap of ρ + (ω + h) and ρ − (ω − h) around ω = 0 shown in Fig. 12(a). While these two DOSs in the unitary limit as well as the weak-coupling side have a relatively large overlap due to the small polaron energy, such an overlap becomes smaller in the strong-coupling side due to the large polaron energy. Physically, the large polaron energy indicates the strong binding of Fermi polarons.
Since there is no single-particle state of majority atoms in the energy range corresponding to low-energy attractive polarons, these polaronic states are irrelevant to spin transport. While the minority DOS at (k F,0 a) −1 = 0.5 is enhanced in the low-energy region (ω < 0) due to the polaron binding effect, that around the energy region where majority DOS ρ + (ω + h) starts to be finite (ω > ∼ 0.2ε F,+ ) is relatively insensitive to the interaction. In addition, the strong attraction makes Z a of the attractive polaron smaller compared with that at unitarity [64], which is expected to reduce the enhancement of spin transport. Therefore, one can find that the non-equilibrium spin transport in highly polarized regime is enhanced by the broadened polaron spectra in the intermediate energy region, and suppressed by the polaron binding effect.
IV. CONCLUSION
In this paper, we elucidated mesoscopic transport properties of spins for strongly interacting Fermi gases connected via a quantum point contact. The tunneling Hamiltonian formalism was used to investigate a steady spin current between two spin-polarized Fermi gases. By employing the linear response theory combined with the diagrammatic approach, the spin current I spin and spin conductance G spin in the normal phase were computed for a wide range of parameters. We found that the emergence of the pseudogap results in a large suppression of spin transport in low-T and low-h regime as shown in Figs. 5-8. On the other hand,the gases become highly polarized for a large spin bias h. In this case, both the attractive polarons at finite momenta and the repulsive branch play significant roles and they lead to the enhancement of I spin compared with the current without self-energy corrections.
As mentioned in previous studies on the spin susceptibility for a strongly interacting Fermi gas [17,[34][35][36][37][39][40][41][42], spin properties are sensitive to the formation of a pseudogap. Mesoscopic spin transport with a small spin bias studied here provides a new probe to experimentally examine the pseudogap phenomenon. We also clarified that the spin current for a large spin bias is affected by excitations including both the attractive polarons at finite momenta and the repulsive branch. At the same time, a set of chemical potentials µ σ,j considered in this article (see Fig. 1) has limitation to examine the properties of polarons in the whole energy range. Such polaronic properties are considered to be accessible in a two terminal system under another choice of µ σ,j with the use of the spin filter, which has been recently realized in an ultracold atom experiment [85]. By assuming µ ↓,L = µ ↓,R µ ↑,L = µ ↑,R +V ↑ and filtering out the majority (σ =↓) component, the fully polarized current for a small bias V ↑ encodes information of attractive polarons at low energy.
Our method can be generalized to fermionic superfluids [86]. Unlike the mass current case, the Josephson current is expected not to contribute to the spin current. Another generalization is the study beyond the linear response theory to discuss the good-contact regime. While such an analysis is generically complicated, the quasiparticle current takes the form of Eq. (2) up to a prefactor in some situation (see Appendix C). Our formalism also predicts the noise of the spin current. Within the linear response theory, the noise at zero frequency is related to I spin by S(ω = 0) = coth(h/T )I spin as derived in Appendix A. We expect such a noise to be accessible in future ultracold atom experiments [87]. This appendix is devoted to the linear response theory to the tunneling amplitude t. For convenience, we introduce C σ ≡ t p,p c † p,σ,R c p ,σ,L andÎ σ (t ) ≡ −dN σ,L (t )/dt . The mass and spin current operators are given byÎ mass =Î ↑ +Î ↓ andÎ spin =Î ↑ −Î ↓ , respectively. The tunneling Hamiltonian in Eq. (1c) is rewritten as H T = σ=↑,↓ (C σ + C † σ ). Using the Heisenberg equation combined with {c p,σ,j , c † p ,σ ,j } = δ pp δ σσ δ jj and {c p,σ,j , c p ,σ ,j } = 0, we obtainÎ σ (t ) = i[N σ,L , H] = −iC σ + iC † σ . First, we review the linear response of the current for the spin-σ component. We follow the procedure in Ref. [78]. Hereafter, · · · denotes the expectation value for a given non-equilibrium state. According to the Kubo formula, I σ ≡ Î σ (t ) for a steady state is given by where · · · eq is a thermal average for the Fermi gases in both reservoirs, and O (A) (t ) ≡ e iAt Oe −iAt has been defined. Operators in the right hand side are in the Heisenberg representation for H 0 = K L + K R + σ,j µ σ,j N σ,j . Using the Baker Campbell Hausdorff formula, we can find C (H0) σ (t ) = e −i∆µσt C (K0) σ (t ), where K 0 = K L +K R and ∆µ σ = µ σL − µ σR . (Henceforth, the shorthand notation C σ (t ) = C (K0) σ (t ) is used only for the operator C σ .) By employing this as well as the expressions ofÎ σ and H T in terms of C σ , Eq. (A1) in the normal phase can be rewritten as where is a retarded correlation function. By using the commutability of K L and K R in Eq. (1) as well as an interrelation between retarded and Matsubara Green's functions, I spin can be rewritten as Eq. (2) in terms of DOSs.
Combining this with Eqs. (A2) and (A6), we obtain In the case of the chemical potentials shown in Fig. 1, we have ∆µ ↑ = −∆µ ↓ = 2h. Therefore, the noise is related to I spin by S(0) = coth(h/T )I spin .
In the end of this appendix, we will comment on the noise of the mass current, which is given by replacing δÎ spin in Eq. (A4) with δÎ mass =Î mass − Î mass . As mentioned above, cross correlation functions between currents with opposite spins vanish in the normal phase. As a result, both mass-and spin-current noises have the same form of Eq. (A6). When a bias is spin independent (µ σ,j = µ j , ∆µ = µ L − µ R = 0), there is no spin current and the noise is related to the mass current by S(0) = coth[∆µ/(2T )]I mass . | 6,441.4 | 2020-02-27T00:00:00.000 | [
"Physics"
] |
Modal Analysis of a Laminar-Flow Airfoil under Buffet Conditions at Re = 500,000
An airfoil undergoing transonic buffet exhibits a complex combination of unsteady shock-wave and boundary-layer phenomena, for which prediction models are deficient. Recent approaches applying computational fluid mechanics methods using turbulence models seem promising, but are still unable to answer some fundamental questions on the detailed buffet mechanism. The present contribution is based on direct numerical simulations of a laminar flow airfoil undergoing transonic buffet at Mach number M = 0.7 and a moderate Reynolds number Re = 500, 000. At an angle of attack α = 4∘, a significant change of the boundary layer stability depending on the aerodynamic load of the airfoil is observed. Besides Kelvin Helmholtz instabilities, a global mode, showing the coupled acoustic and flow-separation dynamics, can be identified, in agreement with literature. These modes are also present in a dynamic mode decomposition (DMD) of the unsteady direct numerical solution. Furthermore, DMD picks up the buffet mode at a Strouhal number of St = 0.12 that agrees with experiments. The reconstruction of the flow fluctuations was found to be more complete and robust with the DMD analysis, compared to the global stability analysis of the mean flow. Raising the angle of attack from α = 3∘ to α = 4∘ leads to an increase in strength of DMD modes corresponding to type C shock motion. An important observation is that, in the present example, transonic buffet is not directly coupled with the shock motion.
aerodynamic lift forces. This aerodynamic instability has been observed in experimental [4,5] and numerical [6] studies of rigid airfoils as well and is known as "transonic buffet". It is generally assumed that the structural response of the wing is triggered by resonance effects after the disturbance amplitude reaches a sufficient magnitude. It is of great interest to be able to define buffet boundaries as precisely as possible in order to fully exploit and potentially extend the safe flight envelope. However, despite large experimental efforts, the self-sustaining mechanism is still not fully understood [1,7]. Transonic buffet is often associated with large amplitude, autonomous shock oscillations, caused by the interactions between shock waves and separated shear layers [1]. While traditional explanations (e.g. acoustic feedback and wave propagation models) have difficulties to directly couple the shock motion with the low-frequency fluctuations in the lift [7], more recent studies describe transonic buffet as a global instability [8]. It is however not quite clear whether the shock motion plays a fundamental active role or is rather a symptom of the periodically accelerating and decelerating flow over the airfoil suction side. Also, "phase-locking" between the shock motion and a global buffet mode is possible. Paladini [9] challenged the importance of the shock motion for transonic buffet, but still suggests that the shock foot plays a major role. We characterise transonic airfoil buffet phenomena here in a more general manner as low-frequency oscillations in the lift coefficient at Strouhal numbers around St = f c/U ∞ ≈ 0.1 (corresponding to typical structural resonance frequencies of aircraft wings) with amplitudes greater than 5%, instead of directly linking it to the shock-oscillation frequency. It should also be noted that buffet on swept wings occurs at higher frequencies [10][11][12] and may be a distinct phenomenon [13][14][15].
In the present contribution direct numerical simulations (DNS) are performed over wingsections at a moderate Reynolds numbers of Re = 500,000 (based on the chord length c) and a Mach number of M = 0.7 considering Dassault Aviation's V2C profile [16]. In the course of the TFAST project, experimental as well as numerical analysis has been carried out on that profile under buffet conditions [17][18][19][20]. There is great interest on global stability analysis of airfoils under buffet conditions (i.e. [6,21,22]), but very little work has been done to date using DNS data. In particular, we wish to explore the usefulness of the dynamic mode decomposition (DMD) procedure [23] to separate and analyse interactions between complex flow phenomena, such as buffet, shock waves, acoustic waves and flow instabilities that are known to exist for this case [16]. Furthermore, it can be useful to compare DMD modes with modes obtained by linear global stability analysis [24]. Based on recent results, we want to extend the investigations for the V2C profile, exploring buffet at angles of attack α = 3 • and α = 4 • (denoted as '3 • case' and '4 • case', respectively), and reduced Reynolds numbers.
After outlining the well-established methodology in the following section, providing references to relevant literature for more details, we review findings and conclusions of previous work on the 4 • case based on Fourier analysis of the pressure and lift histories [16,25]. In Section 4, the sensitivity of convective instabilities to the low-frequency dynamics is analysed, before results from a global stability analysis, and their limitations, are discussed. DMD is applied in Section 5 to extend the study of global modes and the flow dynamics. The sensitivity of DMD modes is then examined considering also DNS results at a decreased angle of attack of α = 3 • . To establish the reduced impact of shock waves on the buffet phenomenon in the present case, an extended Fourier analysis on the 3 • test case is presented in Section 7, before summarising the conclusions in Section 8.
Methodology
All direct numerical simulations in this work were carried out using the high-order fullyparallelised multi-block finite difference in-house code SBLI with details in [26,27]. A recently published paper [16] is concerned with the set-up of airfoil simulations. The dimensionless Navier-Stokes Equations (NSE) are solved by a fourth-order finite difference scheme in space and a low-storage third-order Runge-Kutta scheme in time. The temperature dependency of the dynamic viscosity is modelled by Sutherland's law. Zonal characteristic boundary conditions are applied at the outflow boundaries, while integral characteristic boundary conditions at the remaining domain boundaries avoid wave reflections. In the farfield, an implicit sixth-order filter increases the numerical stability of the simulation. A total variation diminishing (TVD) scheme is used to capture shock waves, but is disabled in boundary layers and near the leading edge. The computational domain is divided into three blocks consisting of one C-block around the airfoil geometry and two H-blocks enclosing the wake-region and outflow. In order to include the blunt trailing edge of the original profile, while maintaining continuous metric terms up to the second order of derivatives, an open-source grid generator was developed and released on GitHub [28]. The reference grid consists of more than one billion points, considering a spanwise domain width of 5%c. The adequacy of the grid resolution is confirmed by a grid-refinement study, based on a spectral error-indicator analysis identifying critical regions in terms of grid-togrid point oscillations [29]. The simulation arrangement, including a grid study confirming also a sufficient domain width of the 4 • buffet case that is further analysed in the present contribution, has already been published in [16].
In order to analyse the linear stability, the flowfield is decomposed into a mean flow with superimposed disturbances. The mean flow is obtained from time-and span-averaged DNS solutions, whereas the disturbances are modelled by a normal-mode ansatz. For local linear stability theory, the compressible Orr-Sommerfeld equations (involving parallel-flow assumptions on the linearised NSE) are solved for a 2D flowfield using the in-house code NoSTRANA [27]. Applying a temporal approach, the solution of an eigenvalue problem provides the temporal growth rate corresponding to an angular frequency for a given set of streamwise and spanwise wave numbers. More details on this methodology, including a full derivation of the equations, flow assumptions, and application to time-and space-averaged baseflows around airfoils are provided by [27] and [30]. For the global stability analysis, the linearised NSE are solved directly, applying a normal-mode ansatz with prescribed spanwise wavenumbers. The large-scale eigenvalue problem is solved in matrix-free mode [31][32][33] using the Implicitly Restarted Arnoldi Method [34] in combination with SBLI as a timestepper. More details on this approach are available in [26].
In order to compare the calculated stability results with the unsteady direct numerical solution of the flowfield, the streaming dynamic mode decomposition (DMD) method is applied [23,24,35]. In contrast to the well-known proper orthogonal decomposition (POD), DMD aims to approximate the non-linear dynamics of an unsteady flow by mapping consecutive snapshots (x t and x t+1 ) of the flowfield by a linear operator (x t+1 = Ax t ) rather than finding an orthonormal basis. The aim is to approximate the non-linear dynamics of a general problem by the best possible linear system which has an exact solution x(t) = ˆ i e ω i t , whereˆ i and ω i are the normalised DMD modes and corresponding complex eigenvalues, respectively. In contrast to the POD method, it does not rely on a model describing the dynamics of coherent structures. DMD for flows consisting of small perturbations of a baseflow can be compared to results from global stability analysis [23]. More details on the implementation of this method are also available in [36].
Unsteady Flow Structures at α = 4 •
The reference simulation, at Re = 500,000 and an angle of attack of α = 4 • , is illustrated in Fig. 1 by means of pressure gradient contours, showing a distinct supersonic region over the upper airfoil surface, bounded by the red sonic line. Two dimensional (2D) simulations already show the formation of strong Kelvin-Helmholtz (KH) vortex structures in the airfoil aft section, initialised by upstream moving pressure waves (also known as Kutta waves) that are caused (in the 2D case) by a von Karman vortex street, which first appear with Strouhal numbers of St ≈ 20 − 30. This suggests a complex cascade mechanism, also involving shock-wave/boundary-layer interaction and the Doppler effect, that allows flow structures of high frequencies to interact with flow phenomena at significantly lower frequencies.
After extruding the 2D solution, a self-sustaining laminar/turbulent boundary-layer transition mechanism sets in on both sides without further artificial excitation of the flowfield. A transition mechanism similar to [37] can be observed, where vortex-stretching of near-wall rib vortices [38] that are lifted up by strong 2D vortices promotes a rapid breakdown to turbulence (inset of Fig. 1). A 2D-like silhouette of the strong vortices can still be observed in the fully turbulent section. Furthermore, those turbulent structures interact with each other as well as the potential flow. In the aft section of the airfoil, strong acoustic radiation can be observed from multiple sources. Black contours in Fig. 1 indicate strong pressure gradients. Approaching the supersonic region, upstream-propagating acoustic waves can be seen to accumulate and form stronger pressure waves (labelled PW in Fig. 1). Eventually, those pressure waves turn into shock waves (labelled SW in Fig. 1) and propagate upstream. Tijdeman [39] distinguishes between three different types of shock behaviour, known as type A, type B, and type C shock motion. Type A shock motion describes a single permanent shock wave oscillating back and forth over large parts of the airfoil, while the shock wave temporarily disappears during the downstream excursion for type B shock motion. Type C shock motion is significantly different from types A and B, as there is no permanent shock wave observed. Instead, there are periodically-generated upstream-propagating shock waves leaving the airfoil at the leading edge that continue moving upstream into the oncoming free stream. All three types have been observed in simulations and experiments. Experimental studies of the same airfoil at higher Reynolds numbers [18] show type A shock motion. The difference between experiments and the present simulations is almost certainly due to Reynolds-number effects [25]. However, confirmation of this would require additional simulations of the same test case at Reynolds numbers of the order of several millions. The present simulations show how acoustic waves circumvent the supersonic region at higher speeds than the upstreampropagating shock waves and introduce additional disturbances into the supersonic region from above. A complex interaction between shock waves, pressure waves, and reflections at the boundary layer is observed. One can also observe Mach-like wave patterns (labelled MW in Fig. 1) on both sides that are caused by acoustic waves travelling upstream within the separated boundary layer (highlighted magenta in Fig. 1). Figure 2a shows the lift coefficient C L as a function of time, highlighting low-lift phases (LLP) and high-lift phases (HLP) in blue and red respectively. The statistical sample size for the HLP and LLP is limited, as the statistics files were only written approximately every two time units (TU) of DNS run time. The time segments for the high-and low-lift phases were selected manually. A distinct low-frequency oscillation in the lift coefficient (up to 12% deviation from the mean value) is observed at a Strouhal number of St = 0.12, which is typical of buffet frequencies [7]. The low-frequency behaviour can also be clearly observed in lift spectra and pressure signals as a function of time at various locations along the surface (including the leading-and trailing-edges) and in the freestream [16]. While high-amplitude lift fluctuations associated with buffet are typically observed at the same frequency as the back-and forth-moving shock waves, there is no obvious correlation in the present case, as upstream-propagating shock waves are generated at significantly higher frequencies (St = 0.4 − 0.7) and leave the airfoil via the leading edge. While studies of the same airfoil at M = 0.7 and significantly higher Reynolds numbers close to experimental conditions (Re = 3 · 10 6 ) do not suggest transonic buffet at α = 4 • , low-frequency oscillations at St ≈ 0.1 are observed for α > 5 • [17,20]. The latter results, comparing delayed detached-eddy simulations (DDES), implicit large-eddy simulations (ILES), and unsteady Reynolds averaged Navier-Stokes (URANS) approaches, show high sensitivity of the low-frequency buffet phenomenon to the modelling methodology applied. This is also reported by an LES study using different grids and methods to reproduce the low-frequency behaviour of the present test case [40].
The black line in Fig. 2b shows the suction-side wall-pressure coefficient C p,w as a function of x, averaged over the full runtime of 25 time units. Phase-averaged C p,w corresponding to low-and high-lift phases (as shown in Fig. 2a) are denoted by the blue dashed and red dash-dotted lines in Fig. 2b, respectively. The main differences are observed at x ≈ 0.7 (region of turbulent breakdown), as the flow is decelerated further downstream by a steeper increase in the surface pressure during HLP compared to LLP. Furthermore, a small plateau is observed at x ≈ 0.4 during HLP, which does not appear in the overall-and LLPaveraged C p,w . The black line in Fig. 2c shows the root-mean-square (rms) of wall-density fluctuations (ρ ) as a function of x, which is calculated at each chord position according to where the total runtime for the present case denotes T tot = 25. For isothermal surfaces, the wall density is directly proportional to the wall pressure so that As already suggested in Fig. 2b, we observe strong fluctuations for 0.2 < x < 0.4 and 0.6 < x < 0.8. To better understand the contribution of the low-frequency phenomenon to the total wall-density fluctuations, the instantaneous wall-density signal is filtered by applying a Fourier band-pass filter. After transforming one-dimensional time signals of ρ at each xposition along the airfoil surface into Fourier space, the power-spectral density is set to zero beyond selected frequency ranges and then transformed back to physical time space. The root-mean-square is again calculated according to Eq. 3 using filtered wall-density fluctuations instead of ρ . The fluctuations corresponding to the low-frequency phenomenon with St < 0.2 (denoted by the blue line in Fig. 2c) are dominant in the fore part of the airfoil (x < 0.4) and at x ≈ 0.7. The contribution of the low-frequency phenomenon to the overall ρ w,rms drops significantly between these regions, where upstream-propagating shock waves are observed in DNS visualisations generated at frequencies in the range 0.4 < St < 0.7. Also in the turbulent part for x > 0.8, fluctuations with St > 0.2 (corresponding to the red line in Fig. 2c) are significantly stronger. While the blue line corresponding to the lowfrequency fluctuations exhibits a distinct peak at x ≈ 0.7, the red line shows a plateau for 0.65 < x < 0.75, where upstream propagating shock waves are formed. Based on Fig. 2c, no direct correlation between fluctuations at low-frequencies for St < 0.2 (corresponding to transonic buffet) and higher frequencies (corresponding to convective structures and shock waves) can be established. It is, however, not possible to rule out a potential interaction between the buffet phenomenon (St = 0.12) and phenomena at significantly higher frequencies (particularly near x = 0.7). Therefore, it is necessary to study the physical phenomena, such as shock waves, acoustic waves, and convective structures more precisely by other means. The full spatio-temporal behaviour of shock-and pressure-waves along the white-dashed line over the airfoil suction side in Fig. 1 is provided by Fig. 3. Figure 3a shows the probability of compression-(red bars) and expansion-waves (blue bars) occurring as a function of the x-location. Cells cut by the sonic line are assigned a value of unity, while remaining cells are denoted by zero. An integration over time suggests how much time shocks spend at a given location. The probability is obtained by dividing the result by the total runtime of 25 time units. A value of unity would imply a steady compression-or expansion-wave occurring at that fixed position for the full time span. As mentioned before, there are no steady or permanently visible shock waves observed at this flow condition. Even though the behaviour of the shock waves seems rather chaotic, there are regions on the airfoil where shock waves tend to spend more time. Considering the red bars, corresponding to compression waves, one can observe increased probability of shock waves occurring at x ≈ 0.42, x ≈ 0.37, just before half chord at x ≈ 0.48, and within the transition region 0.6 < x < 0.775. It is notable that the red line in Fig. 2 showing ρ w,rms for St > 0.2 is rather smooth and does not exhibit local peaks as suggested by Fig. 3a.
The pressure gradient along the gridline at η = 200 (defined by the dashed white line in Fig. 1) is plotted as contours in Fig. 3b as a function of the chordwise location (x) and simulation time (t). Green curves indicate the sonic line (iso-curves for M = 1) and red contours indicate strong adverse pressure gradient (∂p/∂x > 0), while blue contours indicate favourable pressure gradients. Figure 3b also shows strong downstreamconvecting high-frequency pressure fluctuations at the right hand side of the plot (x > 0.7). Near the transition region, a narrow band (centred on the dotted curve in Fig. 3b) of upstream-propagating high-frequency acoustic waves can be observed, which is hidden further downstream by the noise of the large turbulent vortices. These upstream-propagating acoustic waves slow down as they approach the supersonic region, where shock waves are formed. The location of this band is unsteady and varies mainly with the buffet frequency of St = 0.12. Upstream of this band, the flow is mainly supersonic (bounded by the green sonic lines) and shielded from high-frequency oscillations. In the left half of Fig. 3b, pairs of compression-(red) and expansion-waves (blue) are observed, which is also shown by the V-shape of the shocks in the snapshot of Fig. 1. As the waves propagate upstream, the subsonic regions bounded by the green sonic lines between compression-and expansion-waves disappear near the airfoil surface.
The bars in Fig. 3c on the right of b correspond to the number of shock waves counted at each time step. Similarly to Fig. 3a, cells cut by the sonic line are assigned a value of unity, while remaining cells are denoted by zeros. An integration over x gives the values for the red bars in Fig. 3c. The black curve denotes the lift coefficient, importantly showing no direct correlation between the low-frequency buffet phenomenon and the number of apparent shock waves. At least for the present configuration it seems that buffet and shock waves are independent phenomena. The low-lift phases correlate well with the patches of high pressure gradient, at x ≈ 0.9, that were mentioned in connection with Fig. 3b. It can also be observed that strong pressure waves are halted during low-lift phases at x ≈ 0.35, before they continue moving upstream. This continuation of the upstream motion of strong pressure waves sets in during phases where the lift is recovering after a minimum and the flow over the suction side is again accelerated. Further discussion of the unsteady flow phenomena is given in [16,25], and [36].
Local and Global Linear Stability
This section aims to identify linear instabilities and describe them by means of frequencies, growth rates, and spatial coherence. Before attempting global linear stability analysis, a temporal local linear stability theory (LST) approach is applied to time-and span-averaged flowfields of the suction side of a direct numerical simulation of an airfoil at α = 4 • with a total run time of 25 time units. As a consequence of the observed buffet phenomenon, the flowfield and in particular the boundary-and shear-layers change significantly. The contour-plots of Fig. 4a and b show the z-vorticity component (ω z ) of phase-averaged flowfield considering only high-lift phases (HLP) or low-lift phases (LLP) according to Fig. 2a, respectively. Due to the limited number of simulated low-frequency cycles, there are still traces of instantaneous flow features (especially in the freestream) that are not completely averaged out. During low-lift phases, the separation bubble moves upstream and the flow separation becomes more pronounced. The local peak of the blue line in Fig. 2c associated Fig. 4 a Phase-averaged vorticity fields for high-lift phase and b low-lift phase. c Iso-curves for ω z = 50 of high-lift phase and low-lift phase in red and blue, respectively with low-frequency ρ w,rms at x ≈ 0.3 seems to correspond to the intermittent flow separation in that region. Iso-curves in Fig. 4c show a direct comparison of the shear layers for HLP (red) and LLP (blue). For HLP, the flow follows the contour longer, whereas the shear at the lower corner of the blunt trailing edge is significantly reduced. The shear layer along the suction side surface does not change significantly.
Using a mean flow that is averaged over the total run time of the simulation fails to take the periodic variations in lift into account, and it is clear that shear-and boundary-layer characteristics are significantly influenced by the low-frequency oscillations. Therefore the phase averaged flows shown in Fig. 4 are also analysed with respect to linear instabilities. This is justified on the ground that there still exists a wide frequency separation between the buffet mode and the KH modes. Figure 5 shows the temporal growth rate (ω i ) as a function of surface distance s for an angular wave number of ω r = 125 (St ≈ 20), considering only 2D modes (spanwise wavenumber β = 0). The frequency range agrees with Kelvin-Helmholtz instabilities reported by [16] for the present test case. Similar flow structures (St ≈ 25) in a simulation of a high-pressure turbine vane could also be associated with linear instabilities by [30]. The wavy pattern is likely to be due to upstream-moving pressure waves interacting with the boundary layer combined with short time averaging. We can clearly see that the total time average (black curve) underestimates linear instabilities in comparison to high-(red curve) and low-lift (blue curve) phases. Compared to LLP, the boundary layer at HLP shows higher temporal growth rates peaking closer to the trailing edge. Instantaneous snapshots confirm that KH roll-ups form further upstream at LLP. Particularly the LST results at LLP show increased growth rates locally corresponding to regions with high shock probability in Fig. 3a.
We next analyse the global instability of the 2D mean flow with spanwise wavenumber β = 0. Considering the flow around a NACA 0012 airfoil at Re = 200,000 and M = 0.4, Fosas De Pando et al. [41] report a region in the spectrum that is dominated by distinct equally-spaced frequencies (tonal noise) around a maximum peak at St ≈ 7 that corresponds to a stable global mode. An impulse response analysis by [42] showed the vivid interaction between suction and pressure side at that frequency and suggested a feedback mechanism due to pressure waves that are scattered at the trailing edge and form upstream moving acoustic waves. Analysing the time-and span-averaged mean flow over 24 time units, a similar stable mode can be observed in global stability results at St = 5.89 suggesting a growth rate of ω i = −0.019 (negative growth rates denote damping). The divergence field of that global mode in Fig. 6a shows regions of high growth rates in the separation regions on both sides. This global mode also involves upstream-travelling acoustic waves originating at the trailing edge. While those waves can travel along the pressure side without any restrictions, they are slowed down on the suction side approaching the supersonic region. Near the shock wave that bounds the supersonic region in the downstream direction, those waves are compressed as the phase speed decreases to zero. The pressure waves circumventing the supersonic region slide along the sonic line and introduce disturbances that are reflected at the airfoil surface. This mode thus contains a dynamic coupling between separation regions and the trailing edge, via upstream moving pressure waves. The z-vorticity of this eigenmode is shown in Fig. 6b.
Despite the success in extracting the above global mode, in general we found global analysis to be of limited value for the current case. This is partly because the flow is already undergoing moderate buffet and the linearised approach is no longer relevant. In addition, the convergence of global modes was limited by under-resolution of the acoustic field far from the airfoil. As an alternative, we next consider the DMD approach.
Dynamic Mode Decomposition
In this section, a dynamic mode decomposition (DMD) of the test case with 4 • angle of attack is performed considering two sets of instantaneous snapshots over 25 time units, which are in both cases a combination of 2D xyand xz-planes (located within the shear layers) to capture 3D characteristics of the flow, while keeping the required memory low. The following section will consider the sensitivity of the results to the angle of attack. As the modes discussed in this contribution are mainly 2D, visualisations will focus on the xyplane. After discussing the observed DMD modes, a sensitivity analysis using additional data sets with different sample rates and time segments is presented in the next section, so that the reader can better interpret the observations. Figure 7 and growth rates of DMD modes. Due to the limitations in random-access memory, it is necessary to either keep the sampling rate low to compute a larger number of modes (as a POD projection is applied to approximate the mapping matrix A in Eq. 2), or vice-versa. A lower sampling rate, however, is not able to reconstruct high-frequency modes, as shown in Fig. 7. Consistency between data sets with different sampling rates is established in [36] comparing a mode that appears in both sets.
A DMD mode with a similar frequency (St = 6.39) to the global mode (St = 5.89) and decreased damping at St = 6.39 (marked green in the bottom plot of Fig. 7) is shown in (Fig. 8a) shows acoustic waves originating from the trailing edge on both sides of the airfoil, which are also observed in the global mode in Fig. 6. Both the DMD mode and the global mode show traces of shock waves near the half chord position. The DMD mode in Fig. 8a, however, is very noisy within the shear layers and wake, so that the structures shown by the global mode at x ≈ 0.8 on the suction side and at x ≈ 0.9 on the pressure side are hidden. These structures can be better seen in the velocity and density shown in Fig. 8b and c. After having identified a DMD mode, which is similar to a global mode, we extend the study to other flow features. Selected eigenfunctions of the modes marked red in Fig. 7 are shown in Fig. 9, plotting contours of density (left column) and streamwise velocity component (right column). The top row shows a mode at St = 19.3, picking up the KH rollups, which are associated with convective linear instabilities. The shape of the eigenmode is reminiscent of structures that are observed in movie snapshots like Fig. 1. Figure 7 shows several modes at frequencies between 15 < St < 25 with similar or even lower damping Fig. 9 Selected DMD eigenmodes at different frequencies. The left and right columns show the density (ρ) and streamwise velocity (u) field, respectively rates with similar shapes of eigenmodes. Considering the results of the local linear stability analysis, the frequency of the KH instabilities is expected to vary significantly due to the low-frequency dynamics of the flowfield.
There are big turbulent vortices observed downstream of the laminar/turbulent transition region with Strouhal numbers around St = 1.8. These modes can also be found in the DMD spectrum. The density field in the second row of Fig. 9 shows strong oscillations in the aft section of the airfoil and in the wake. In addition, upstream-moving pressure waves, originating at the trailing edge can be observed. A phase shift of the oscillations within the upper-side shear layer and the freestream can be seen in the velocity field of that mode.
Besides the low-frequency peak at St = 0.12 (the last row of Fig. 9) corresponding to the buffet phenomenon, distinct low-frequency modes at St = 0.6 and St = 0.45 are shown in the third and fourth rows of Fig. 9, respectively. The variation of the spacing between green sonic lines at x ≈ 0.6 in Fig. 3b Fig. 9 are chosen as representative modes for the shock motion, as they have high normalised amplitudes and low damping rates in the DMD spectrum (marked red on Fig. 7). Furthermore, high power-spectral density is also observed in the Fourier spectra of surface pressure probes at these frequencies [25]. Both eigenmodes look similar, but have a clear phase shift. These modes seem to be strongly coupled with the shock dynamics and fluctuations in the wake. Their eigenmodes are mainly active on the suction side, but also include acoustic waves moving over the pressure side and also comprise acoustic waves circumventing the supersonic region. These modes can be interpreted in the light of other observations of higher frequencies on airfoils. Performing a large-eddy simulation of a supercritical laminar airfoil (OALT25) at α = 4 • and M = 0.735, but at a higher Reynolds number of Re = 3 · 10 6 , Dandois et al. [43] reported shock motion (a permanent back and forth moving shock wave) at significantly higher Strouhal numbers (St = 1.2) compared to typical buffet frequencies, and linked it to a breathing phenomenon of the separation bubble associated with downstream convecting vortices. Similar phenomena were reported by Memmolo et al. [17] analysing the V2C airfoil at α = 7 • and high Reynolds numbers using large-eddy simulation. In a URANS parameter study exploring the buffet domain varying the angle of attack (3 ≤ α ≤ 9.5) and Mach number (0.55 < M < 0.75), an interesting coexistence between type A and type C shock motion was reported by Giannelis et al. [44]. The disturbances associated with the type C shock motion emerged from downstream-convecting recirculation pockets that produced oscillations at the trailing edge at significantly higher frequencies compared to the low-frequency lift oscillation associated with buffet. In the present case, an acoustic propagation model, originally proposed by Lee [7] to predict transonic buffet, is able to approximate the frequency range associated with the periodic shock motion, but not the low-frequency buffet phenomenon [36].
The mode with St = 0.12 (in the last row of Fig. 9) agrees with typical transonic buffet frequencies in the literature and can also be extracted from the lift coefficient over time.
In particular, the density fluctuations in the rear part of the airfoil agree well with observations by Fukushima and Kawai [45]. Strong fluctuations of the modes associated with the shock motion at x ≈ 0.6 are located slightly upstream of high amplitudes in the buffet mode (St = 0.12) at x ≈ 0.7. High amplitudes are also observed in supersonic regions around x = 0.4, were Sartor et al. [46] reported high receptivity of the global buffet mode along characteristic lines for a OAT15A airfoil at buffet conditions. In the velocity field, the suction and pressure sides are clearly separated by a phase shift, indicating an opposed streamwise oscillation of the shear layers. Furthermore, there seems to be a phase shift between the separation regions and the shear layers. In particular, the velocity field of the Fig. 10 Lift coefficient of test cases at 4 • and 3 • angles of attack, corresponding to black curves and red curves, respectively. The lighter coloured lines were disregarded from the statistics DMD buffet mode is qualitatively very similar to the unstable global buffet mode reported by [47]. Density fluctuations of the modes with St < 1 are not only observed in the streamwise direction, but also in the wall-normal direction.
Fluctuations near the trailing edge in DMD modes corresponding to shock motion (St = 0.45 and St = 0.6) are reminiscent of the buffet mode at St = 0.12. The difficulties in fully separating shock motion from the buffet phenomenon using DMD can be either due to the sensitivity to the selected data sets, or to large lift oscillations at well-established buffet conditions. Therefore, the robustness of DMD to sample rates and sampling time needs to be studied, before analysing a test case closer to onset conditions at α = 3 • in order to establish the (minor) impact of shock motion on transonic buffet in the following sections.
Sensitivity Study of DMD Results
Dynamic mode decomposition has been widely used to study the flow dynamics over airfoils (e.g. [48][49][50]). Comparing DMD with POD methods, [51] found DMD favourable to study transonic buffet. In this section additional data sets of the simulation at α = 4 • are analysed using DMD. Furthermore, a simulation is performed at a decreased angle of α = 3 • , using the same DNS set-up as described above to study the sensitivity of lowfrequency phenomena (observed at St < 1) to the angle of attack. Figure 10 shows the lift coefficient as a function of time for the 4 • case (grey and black lines) and 3 • case (red lines). The 3 • case was restarted from a solution of the 4 • case at t = 13 and underwent a transient between t = 13−32 (denoted by the light red-coloured line), before reaching the developed quasi-periodic flow (denoted by the dark red line). For the mode analysis it is useful to perform a sensitivity analysis of the dominant DMD modes to the sampling period. In order to save random-access memory, it would be also interesting to find out whether 2D xy-slices alone are sufficient to accurately capture the DMD modes corresponding to low-frequency phenomena. Therefore, three data sets of the 4 • case are analysed and summarised in the first three rows of Table 1, before comparing the DMD modes of data sets at different angles of attack (last two rows in Table 1).
Data set 4a was discussed in the previous section in connection with the top plot of Fig. 7 and the bottom three DMD modes in Fig. 9. Data set 4b contains the same number of snapshots (250) and sampling rate (sampling every 0.1 time units), but uses only xyslices. Both data sets 4a and 4b consider the full runtime of 24 time units (corresponding to the grey and black lines in Fig. 10) sampling 10 snapshots per time unit. Data set 4c has a slightly lower sampling rate of 8.33 snapshots per time unit (sampling every 0.12 time units) using 140 xy-slices over 16 time units (corresponding to the black line in Fig. 10). Data set 3a comprises data from the test case at a lower angle of attack and can be compared to data set 4c. Figure 11a shows the normalised amplitude of DMD modes for the case at α = 4 • , where data set 4a is represented by black rectangles, while blue and red triangular symbols correspond to data sets 4b and 4c, respectively. Perfect agreement between all three data sets can be observed for the buffet mode at St = 0.12 (outlined red). Also the DMD mode at St = 0.6 seems to be consistent, while the amplitude slightly varies (modes are outlined green). For data set 4a, we find a DMD mode with a high amplitude at St = 0.45. Looking at the DMD modes within the magenta ellipse in Fig. 11a, data set 4b again shows good agreement, whereas data set 4c with a shorter sampling time suggests modes at slightly higher and lower frequencies (St = 0.42 and St = 0.61, respectively). This is due to the variation of time scales corresponding to the shock motion, which was also observed in Fig. 3b. DMD modes found in the range of St = 0.4 − 0.6 are, however, qualitatively similar and can be clearly linked to the shock behaviour. Data set 4c gives a DMD mode with increased amplitude at St = 0.24 (outlined by the blue ellipse in Fig. 11a). This DMD mode is thus a super-harmonic of the buffet mode (St = 0.12). Varying the sampling rate, [51] also reported low sensitivity of low-frequency DMD modes (St < 1) for sampling rates greater than 3 snapshots per time unit (corresponding to 24 snapshots per buffet cycle).
After having shown the robustness of DMD modes to the selection of data sets, Fig. 11b shows the normalised amplitudes comparing DMD modes at different angles of attack considering a sampling time of 16 time units. Similar to Fig. 11a, red symbols correspond to Fig. 11 Normalised amplitude of low-frequency DMD modes (St < 1) considering a three data sets of the 4 • case (first three rows in table 1), and b data sets of the 4 • and 3 • test cases (last two rows in table 1). The modes corresponding to the filled symbols are shown in Fig. 12 α = 4 • and black circles to α = 3 • . The filled symbols are the modes of interest. Again, we can see good agreement for the buffet mode St ≈ 0.12. At lower angles of attack, the buffet frequency seems to slightly increase, while the amplitude decreases. Also the lift fluctuations in Fig. 10 for the 3 • case are smaller than for the 4 • case. Even though the frequency of well-established buffet oscillations often increases with the angle of attack [1], some experimental studies of laminar-flow airfoils near buffet-onset conditions [52] show a similar trend as for the present test case. For the α = 4 • simulation, this can also be due to the short sampling time. A significant decrease in normalised amplitude is observed for the DMD modes associated with the shock motion and there is no significant peak observed in the DMD spectrum at α = 3 • in the range of St = 0.4 − 1. Figure 12 shows density fluctuations of the DMD modes corresponding to the black filled circles and red filled triangles of Fig. 11b A comparison of DMD modes at two different angles of attack (α = 3 • and α = 4 • ) suggests that modes corresponding to shock motion (St = 0.4−0.7) have increased amplitudes for higher angles of attack. For the α = 4 • case, these modes contain features reminiscent of the buffet mode at St = 0.12 (marked by x-symbols in Fig. 12). These strong fluctuations near the trailing edge are significantly less pronounced at α = 3 • . This suggests a potential coupling or phase-locking of shock motion and transonic buffet at higher angles of attack, which would lead to the traditional observations of a single shock wave oscillating back and forth at a distinct buffet frequency corresponding to large-amplitude lift oscillations [18-20, 53, 54]. In the present cases (in particular the 3 • case), there is no clear evidence on an interaction between shock waves and low-frequency lift oscillations. Instead, the shock motion appears as a consequence to acceleration and deceleration of the flow due to the buffet mode, as regions of high amplitudes in the DMD modes at St = 0.4 − 0.7 are located between regions of strong fluctuation in the buffet modes (marked by dotted lines and xsymbols in the bottom plots of Fig. 12). The following section aims to confirm the minor effect of the shock motion on the surface pressure.
Fourier Analysis of the Wall Density
After having gained a good impression of the global dynamics in the xy-plane, we now perform Fourier analysis along representative gridlines with z = 0 to study acoustic phenomena outside the boundary layers and, in particular, their foot-print on the wall density (which is proportional to the wall pressure for isothermal boundary conditions). We focus on the 3 • case, as the amplitudes of large-scale fluctuations in the lift coefficient in the range of 27 < t < 52 (denoted by the dark red line in Fig, 10) are more regular compared to the 4 • case. An instantaneous snapshot in Fig. 13a at t = 37.9 shows contours of the streamwise density gradient during the high-lift phase with the sonic line in green. Flow quantities are monitored on the airfoil surface, and along the dashed and dash-dotted lines, to generate similar space/time diagrams as in Fig. 3b for the 4 • case. Contours of the streamwise density gradient are shown in Fig. 13b as a function of x and time, where the horizontal dashed black line corresponds to the dashed black curve in Fig. 13a. Similar plots are shown in Fig. 13c and d showing the density gradient as a function of x and time along the surface and along the dash-dotted line outside the supersonic region in Fig. 13a, respectively. The green lines in Fig. 13b-d correspond to the sonic lines as a function of x and time along the dashed black line in Fig. 13a and represent the shock dynamics. While the wall-density gradient in Fig. 13c represents a superposition of the buffet phenomenon, downstream-convecting turbulent structures, and traces of acoustic phenomena, Fig. 13d comprises mainly acoustic pressure waves circumventing the supersonic region (sketched by the black-dotted curves), which propagate upstream significantly faster than the shock waves indicated by the green lines.
To get a clearer picture of the dynamics of the wall density, the same Fourier low-pass filter technique from Section 3 is applied to wall-density fluctuations for the 4 • case. Filtered space/time diagrams at the surface are shown in Fig. 14 In Fig. 14a, large amplitudes in the low-frequency wall-density fluctuations on the suction side are observed at x ≈ 0.4 and at x ≈ 0.8, which agree with observations from the DMD analysis. We observe significantly lower amplitudes in regions near flow separation (0.4 < x < 0.7). Before and after that region, we can observe respectively the upstreamand downstream-propagation of peaks and troughs. The propagation speeds, however, are significantly slower than either the convection or acoustic velocities, which can be observed in Fig. 13b and d (sketched by dotted curves). While the fluctuations on both sides are in phase near the trailing edge in Fig. 14a, a clear phase shift can be observed in transition regions on the pressure and suction sides at x ≈ 0.8 and x ≈ 0.7, respectively. The surface-density fluctuations upstream of transition regions are in phase with the shear-layer dynamics, whereas the strong fluctuations near the trailing edge are 180 • phase shifted, as shown in the DMD mode at St = 0.15 in Fig. 12. While the low-frequency fluctuations for x > 0.5 on the pressure side seem to be stationary, for x < 0.5 they propagate upstream at similar speeds as acoustic waves. A comparison of low-frequency oscillations on the suction-side surface with the shock motion, illustrated by the green lines in Fig. 14a, suggests that the global buffet mode sets the speed of upstream-propagating shock waves. During high-density phases (red contours), upstream-propagating shock waves are able to Fig. 13a. The dotted curve in (a) suggests the potential path of a single shock wave maintain subsonic regions between shock-and reflection-waves further upstream than during low-density phases (blue contours). The dotted curve in Fig. 14a sketches the potential path of a single shock wave oscillating between strong fluctuations at x ≈ 0.4 and x ≈ 0.8, always moving from alternating high-to low-density peaks. Whether this appears at higher Reynolds numbers requires further investigations. Figure 14b, showing phenomena in the range of 0.2 < St < 1, suggests that the surface density is not significantly affected by shock waves. Even though upstream-propagating shock waves are observed in this frequency range, the green lines associated with the shock motion do not align with the contours of the wall-density gradient. Instead, we can observe traces of upstream-propagating acoustic waves circumventing the supersonic region at significantly higher speeds compared to the upstream-moving shock waves (as was shown in Fig. 13d). Towards the rear of the suction side, downstream-convecting structures appear in this frequency range, which are associated with large turbulent structures. These structures generate new acoustic waves when they arrive at the trailing edge.
The Fourier analysis in this section suggests that the shock motion is modified by the global buffet mode, as the irregular shock pattern (corresponding to St = 0.4 − 0.7) follows large-scale surface-density fluctuations associated with the St = 0.12 phenomenon.
Conclusion
DNS data has firstly been analysed in terms of local and global linear stability in order to investigate the transonic buffet mechanism of a narrow wing section at α = 4 • and Mach and Reynolds numbers of M = 0.7 and Re = 500,000, respectively. Local linear stability theory associates KH roll-ups seen in the DNS with convective linear instabilities, as has been previously observed for a high-pressure turbine vane at similar freestream conditions [30]. The shear layer on the suction side shows significantly different characteristics during high-and low-lift phases. The analysis of the time-and span-averaged flowfield underestimates the growth rates of those instabilities compared to phase-averaged flowfields for high-lift and low-lift phases. During the high-lift phase, the unstable region in the boundary layer is further upstream, with higher growth rates at higher frequencies, compared to the low-lift phase.
Global stability analysis captures a tonal mode at St = 5.89 that has been previously reported in literature with respect to the coupled dynamics of separation bubbles [42], but other modes fail to converge. Dynamic mode decomposition is able to capture the KH instabilities as well as the global mode at St = 5.89. Furthermore, the DMD shows an eigenmode at St = 1.87 that is associated with large turbulent vortices in the suction-side aft section of the airfoil. Those roll-ups seem to be coupled with the shock region and large-scale structures in the wake that are also observed in instantaneous visualisations. A phase shift in the streamwise velocity field of that mode occurs on the upper side between the wake and freestream. Modes at lower frequencies with 0.3 < St < 1 have similar structures with high amplitudes in similar regions and correspond to flapping of the separated shear layers and involve the shock motion as well as acoustic waves circumventing the supersonic region.
At a lower angle of attack, these spectral intermediate peaks are less pronounced. A DMD mode at St = 0.44 mainly describes the shock motion, whereas DMD modes at higher frequencies again comprise shock motion as well as acoustic waves. The existence of a bifurcation in the shock dynamics was already suggested by Giannelis et al. [44], where at certain flow conditions upstream-propagating pressure waves are periodically generated at the trailing edge at significantly higher frequencies than the co-existing permanent shock wave oscillating at a typical frequency associated with buffet (St ≈ 0.1). In the present case, no permanent shock wave is apparent so that only type C shock motion is observed, similar to experiments by McDevitt at al. [55]. The DMD modes associated with buffet at St = 0.12 − 0.15 show high fluctuations in the rear part of the airfoil and around x ≈ 0.4, but do not seem to be directly coupled with the shock motion. In contrast to the modes at high frequencies, the buffet DMD mode does not change significantly when decreasing the angle of attack.
The DMD buffet modes are reminiscent of the global modes reported by Crouch et al. [8] and Sartor et al. [46]. Propagation speeds of low-frequency structures do not correspond to acoustic speeds, suggesting that Lee's acoustic feedback mechanism of buffet [7] is not active here. Instead it seems that the shock motion is modulated by the low-frequency buffet. The current observations suggest that transonic buffet does not rely on an interaction with shock waves, but does not exclude interactions and phase-locking between low-frequency phenomena and shock waves. Simulations at higher Reynolds numbers or higher angles of attack would be useful to see if there is a change to a single shock wave (visible at all times) oscillating between high-and low-pressure regions as sketched by the dotted curve in Fig. 14a.
Based on the conclusions of this work, it may be necessary to revisit the definition of transonic buffet to distinguish between transonic "shock buffet", where a single shock wave oscillates at a frequency that agrees with low-frequency oscillations in the lift signal, and a form of "incipient buffet", where the shock motion is not necessarily correlated with the dominant oscillations in the lift signal.
Open Access This article is distributed under the terms of the Creative Commons Attribution 4.0 International License (http://creativecommons.org/licenses/by/4.0/), which permits unrestricted use, distribution, and reproduction in any medium, provided you give appropriate credit to the original author(s) and the source, provide a link to the Creative Commons license, and indicate if changes were made. | 11,806.2 | 2019-12-03T00:00:00.000 | [
"Physics",
"Engineering"
] |
Gene Coexpression Network Analysis as a Source of Functional Annotation for Rice Genes
With the existence of large publicly available plant gene expression data sets, many groups have undertaken data analyses to construct gene coexpression networks and functionally annotate genes. Often, a large compendium of unrelated or condition-independent expression data is used to construct gene networks. Condition-dependent expression experiments consisting of well-defined conditions/treatments have also been used to create coexpression networks to help examine particular biological processes. Gene networks derived from either condition-dependent or condition-independent data can be difficult to interpret if a large number of genes and connections are present. However, algorithms exist to identify modules of highly connected and biologically relevant genes within coexpression networks. In this study, we have used publicly available rice (Oryza sativa) gene expression data to create gene coexpression networks using both condition-dependent and condition-independent data and have identified gene modules within these networks using the Weighted Gene Coexpression Network Analysis method. We compared the number of genes assigned to modules and the biological interpretability of gene coexpression modules to assess the utility of condition-dependent and condition-independent gene coexpression networks. For the purpose of providing functional annotation to rice genes, we found that gene modules identified by coexpression analysis of condition-dependent gene expression experiments to be more useful than gene modules identified by analysis of a condition-independent data set. We have incorporated our results into the MSU Rice Genome Annotation Project database as additional expression-based annotation for 13,537 genes, 2,980 of which lack a functional annotation description. These results provide two new types of functional annotation for our database. Genes in modules are now associated with groups of genes that constitute a collective functional annotation of those modules. Additionally, the expression patterns of genes across the treatments/conditions of an expression experiment comprise a second form of useful annotation.
Introduction
The importance of large-scale gene expression analysis in understanding gene function became apparent with the first report of genome-wide transcript expression profiling with DNA microarrays [1]. This led to the use of coexpression analyses to measure not only the physiological state of cells but also to characterize genes with no known function [2]. As more gene expression data sets became available, data from multiple experiments were combined into single analyses to functionally annotate genes based on the conditions under which they are expressed and their correlation to genes with similar expression patterns [3,4]. In plants, numerous projects perform large-scale gene expression analyses in which coexpression networks are created. Several of these combine results from individual experiments and utilize Pearson correlation coefficients between all gene pairs [5,6,7,8,9,10,11] while others incorporate multiple types of data including gene transcript levels, protein-protein interactions, metabolite profiles, and predicted conserved gene interactions [6,12,13,14].
A number of publicly available gene coexpression network databases have been constructed that allow researchers to query pre-constructed gene networks with a target gene(s). These databases permit the identification of correlated gene partners and visualization of a graphical display of coexpression networks with user-specified cutoff criteria including specific experiments or conditions upon which the correlation calculation is performed [5,6,7,8,11]. One confounding problem with current analysis and display methods is that coexpression networks can be very complex thereby making interpretation difficult. Although the selection of a correlation value cutoff can simplify a network by reducing the number of edges, the understanding of gene networks is still problematic [15,16]. Due to the complexity of gene coexpression networks, various methods have been used to find the most informative relationships within correlation networks [17,18,19,20,21,22,23,24].
Several research groups have identified subsets of highly correlated genes within large gene coexpression networks in Arabidopsis thaliana and rice (Oryza sativa) [14,17,18,19,21,22,25,26]. Using various algorithms, these reports examine gene coexpression networks to identify subsets of genes that are more highly connected and highly correlated to each other than they are to other genes in the network. These subnetworks of genes are referred to as modules. Genes within such modules have been shown to be enriched for particular Gene Ontology (GO) categories [17,18,19,22], and relationships depicted by gene modules are congruent with expected gene pathways [18,19,22]. Additionally, hypotheses formulated from gene coexpression modules for particular genetic pathways related to seed embyro development, chlorophyll degradation, organ development and lectin receptor kinase inhibition of seed germination have been substantiated by downstream laboratory experiments [18,27,28,29].
Methods for analyzing genome-wide expression data are either condition-dependent or condition-independent depending on the selection of input data. Condition-dependent data consist of planned treatments/conditions that are designed to record transcript responses to specific physiological states. In contrast, condition-independent data are a compilation of unrelated treatments/conditions that are not designed to provide insight to a particular biological response. Most large-scale plant gene coexpression resources utilize condition-independent analyses that rely upon large compendia of gene expression data sets from independent sources [6,7,8,9,10,13,15,17,18,19,21,22,25,26]. Such analyses are convenient because they make use of the maximal available data. However, there are potential problems with condition-independent analyses as it has been demonstrated that gene coexpression analysis with too many microarray samples can result in the loss of information [30]. Difficulty in interpreting the biological meaning of correlations in complex conditionindependent data sets is a second problem with this analysis strategy. In contrast, condition-dependent analyses typically utilize a smaller, defined set of treatments or conditions that have been chosen to test a particular hypothesis or offer insight into a specific physiological condition [15,16]. Nonetheless, both conditionindependent and condition-dependent gene coexpression studies have utility. Analyses from large condition-independent data sets are likely to identify highly conserved core gene networks while smaller condition-dependent experiments offer the opportunity to recognize more narrowly defined correlations.
In this study, we have adopted a condition-dependent approach and have separately analyzed fifteen rice gene expression data sets based on the Affymetrix Gene-Chip Rice Genome Array using Weighted Gene Correlation Network Analysis (WGCNA), a network analysis method that has been widely used to identify biologically meaningful gene modules in a variety of organisms [24,31,32,33,34,35,36]. Additionally, we created a conditionindependent data set from the same fifteen rice gene expression experiments and identified gene modules from the combined data. A comparison of the results from the two analyses suggests that while both have utility, the data analysis from individual experiments facilitates biological interpretation and is less likely to obscure uncommon but potentially informative gene coexpression modules than the combined data set. Using the conditiondependent results, we have supplemented the annotation of rice genes as 17,298 of the 40,829 protein coding genes in the MSU Rice Genome Annotation Project lack assigned functional annotation [37]. These results provide two important types of annotation. Genes included in these analyses are now associated with expression patterns across defined treatments/conditions. Additionally, genes that have been assigned to coexpression modules can be considered in the context of all other genes that are found within the same module. Both module membership and individual gene expression patterns have been incorporated as part of the annotation in the MSU Rice Genome Annotation Project database (http://rice.plantbiology.msu.edu) [37].
Datasets Used in This Study
Publicly available rice gene expression data were downloaded from the National Center for Biotechnology Information Gene Expression Omnibus (NCBI GEO) and European Bioinformatics Institute (EBI) ArrayExpress [38,39] in February 2010. Only data that had been generated using the Affymetrix Rice GeneChip were considered for analysis. In total, fifteen data sets were chosen for analysis in this study representing 440 arrays (Tables 1, S1, S2). The experimental conditions used to generate the data sets included biotic and abiotic stresses, cytokinin treatment, gibberellin signalling pathway mutant analysis, an extensive tissue atlas, seed germination time courses, an inflorescence and seed developmental series, and photoperiod/thermoperiod time courses [40,41,42,43,44,45,46,47,48,49,50]. Not all samples or treatments/conditions for each data set were included in the analyses. In a few experiments, some treatments/conditions were excluded in order to simplify the interpretation of the results. For example, only expression data for a single rice cultivar, Minghui 63, were included in the analysis of the GSE19024 tissue atlas. Also, root and leaf samples were not essential for the GSE6893 inflorescence and seed developmental series, and root and leaf samples were removed from the dataset. Some individual chips were also excluded after quality analysis (see Materials and Methods), and in two cases, this resulted in all replicates for a single treatment being discarded: shoot 2Fe+P from GSE17245 and LL LDHC 124 hrs from E-MEXP-2506. Descriptions of the chips that were analyzed for each experiment in this study as well as the number of arrays and samples/treatments per experiment are provided in Tables S1 and S2.
Data from each experiment were analyzed individually or as a single combined data set using the WGCNA method [24]. The goals of the analyses were to identify modules of highly coexpressed genes using both methods (condition-dependent and condition-independent) and then to select the method with the most informative results for supplemental rice gene annotation. For both methods, normalized trend plots were generated for all gene modules. WGCNA analyses were assessed by the number of modules identified, the similarity of expression values for the genes within a module, and the biological interpretability of the expression patterns of the genes within modules. Although relaxation of WGCNA-required parameters would have resulted in additional genes being assigned to modules, this would have reduced the overall correlation of the genes in each module (see Materials and Methods, Table 1).
Coexpression analyses from individual, conditiondependent experiments
Following coefficient of variation (CV) filtering of the conditiondependent experiments, a total of 13,537 genes were retained for gene coexpression analysis in at least one experiment (range 672 to 7,478; Table S3). From all 15 experiments, 71 coexpression modules were identified containing 12,328 non-redundant genes ( Table 2, Figures 1, 2, S1, S2, S3, S4, S5, S6, S7, S8, S9, S10, S11, S12, S13). The remaining 1,209 genes that passed CV filtering were not assigned to any coexpression module. The number of modules identified within an experiment varied from two to nine, and the number of genes assigned to all modules within a single experiment ranged from 567 to 4,566. Modules contained between 40 and 3,574 genes with an average module size of 405 genes. The majority of genes assigned to coexpression modules have functional annotation, but nearly one fifth (2,908) of all genes assigned to modules lack functional annotation. Transposable element (TE) related loci were included in the gene sets for these analyses, but overall, only 406 of the genes assigned to modules were TE-related (Table 2), consistent with their reduced levels of expression. While a gene can be present in only one module from a single experiment, many genes were found in multiple modules from different experiments ( Table 3). In fact, most genes that had been assigned to modules were found in modules from two or more experiments, and one gene, LOC_Os11g31540, a BRASSI-NOSTEROID INSENSITIVE 1-associated receptor kinase 1 precursor, was found in modules from 12 different experiments ( Table 3).
The gene coexpression modules identified from the panicle and seed developmental series (GSE6893, [42]) are illustrative of the results that can be obtained using WGCNA analysis with coexpression data. Expression values from a total of 4,231 genes were analyzed from this experiment (Table S3). Eight modules were identified, and the number of genes per module ranged from 104 to 725 with 1,223 genes not assigned to any module. The expression patterns for each module are distinctive ( Figure 1). Some modules coincide with very specific periods of growth such as anthesis ( Figure 1H), middle seed development ( Figure 1D) or late panicle maturation ( Figure 1E). Two modules show gene expression levels that are elevated during both panicle and seed development ( Figures 1A, 1C). Three modules contain genes that are both positively and negatively correlated and that have expression levels that are alternately high and low in panicles and seeds ( Figures 1B, 1F, 1G).
Gene modules obtained by analysis of expression data from a pathogen response experiment (GSE10373) are shown in Figure 2 [43]. This time course experiment was performed on two rice genotypes, Nipponbare and IAC165, after two treatment conditions, mock inoculation and infection with the parasitic weed Striga hermonthica. Because the samples were all derived from the same tissue type (roots), fewer genes (672) passed the CV filter relative to the developmental time course that contained a variety of tissue types ( Figure 1). The genes were split into three modules ranging in size from 52 to 351 (Table S3) that display either genotype by treatment responses (Figures 2A and 2B) or genotype specific expression ( Figure 2C).
Enrichment analysis was performed to identify genes containing particular Pfam domains that are over-represented in these coexpression modules (Tables 4, S4). Statistically significant enrichment was observed in modules from all 15 experiments analyzed. A total of 61 modules were found to have enrichment of genes with at least one Pfam domain, and 114 Pfam domains were enriched in at least one module. A number of modules had enrichment of Pfam domains consistent with the assayed biology. For example, the GSE6893-blue module contains genes that are expressed during late seed development ( Figure 1B) and enrichment of genes with seed-related cupin, protease inhibitor/ seed storage/LTP family and starch synthase catalytic Pfam domains was evident (Table S4) [51,52,53]. Also, the GSE10373blue, GSE16793-blue and GSE18361-blue modules have higher than expected numbers of genes with terpene synthase, WRKY DNA binding and chitinase domains, all domains that are found in genes that are known to be responsive to biotic stresses (Tables S4, S5) [54,55,56,57].
Coexpression analyses from combined, conditionindependent experiments
A condition-independent data set was constructed by combining all data from the 15 condition-dependent experiments used above and performing coexpression analysis with WGCNA. After CV filtering 17,320 genes were used for gene module identification using WGCNA. Only 15 modules containing 10,077 genes were Identifiers for data are from either NCBI GEO or EBI ArrayExpress. 2 Coefficient of variation cutoff used to filter averaged and normalized gene expression data. 3 Beta and tree cut parameters used during WGCNA analysis of expression data. 4 Only shoot apical meristem, developing panicle and developing seed samples were used for this analysis. 5 Shoot 2Fe+P samples were removed after chip QC analysis. 6 Only data from Minghui 63 were analyzed. Expression data from Zhenshan 97 were excluded from analysis. Callus tissue samples were not included in the analysis. 7 The LL-LDHC-124 hrs sample was excluded from analysis after chip QC analysis. doi:10.1371/journal.pone.0022196.t001 identified from the combined data set (Tables 2, S6). Those modules varied in size from 40 to 3,740 genes and had an average size of 671 genes. There were 7,481 non-TE related genes with functional annotation and 2,403 genes with no functional annotation assigned to modules. Enrichment analysis was also performed to identify Pfam domains that were over-represented in genes from the condition-independent coexpression modules. A total of 14 modules had enrichment of a total of 209 Pfam domains (Table S7).
In combination, the condition-dependent and condition-independent analyses included 18,598 genes, of which 15,336 were assigned to at least one module from at least one analysis. Of the 12,259 genes common to both types of analysis, 11,204 were assigned to modules from the condition-dependent experiments, but only 7,480 were found in condition-independent modules. Modules from both the condition-dependent and conditionindependent analyses contained a common subset of 7,069 genes. There were 5,259 genes found in at least one condition-dependent module that were not assigned to any modules from the conditionindependent analysis and 3,008 genes found in a conditionindependent module that were not found in any conditiondependent modules ( Figure 3).
Fewer genes were assigned to gene coexpression modules from condition-independent compared to condition-dependent analyses, and there were fewer modules identified from the conditionindependent analysis ( Table 2). An examination of the trend plots of the condition-independent gene modules shows that some of the patterns observed in condition-dependent gene modules can be observed in condition-independent modules (e.g., Figure S9B vs. Figure S14A; Figure S9B vs. Figure S14B; Figure S13F vs. Figure S14C). Additionally, some condition-independent modules have similar gene expression patterns across a subset of conditions. Figures S14A and S14B show gene expression patterns from the green-yellow and pink modules from the condition-independent analysis, and these modules have similar patterns of gene expression across numerous samples. However, some striking expression patterns from condition-dependent modules are not easily identified in any condition-independent modules such as the anti-correlated circadian cycles in Figures S13E and S13I or the infection response expression in Figure S6A; these expression patterns may be obscured within a densely populated conditionindependent module. A figure containing all gene expression trend plots for each condition-independent gene module can be downloaded from the MSU Rice Genome Annotation FTP site (ftp://ftp.plantbiology.msu.edu/pub/data/rice_gene_assoc/Figure _condition_independent_modules.pdf).
A comparison was made to identify the overlap in genes between modules from the two strategies (Table S8). Often, a high proportion of genes from individual experiment modules were assigned to a gene coexpression module from the conditionindependent analysis. This is not absolute as fewer than half of the genes from some condition-dependent modules were present in the condition-independent modules. In a few cases, the majority of genes from a condition-dependent module were almost entirely contained within a single condition-independent module. However, the more common occurrence was for genes from a single condition-dependent module to be distributed between a subset of condition-independent modules, and this was the case for the modules described above in Figures S9B, S14A, S14B, which represent the GSE19024-brown module and the conditionindependent green yellow and pink modules (Table S8).
Improvement of rice gene annotation via coexpression analyses
We incorporated the results from the analyses of individual condition-dependent experiments into the MSU Rice Genome Annotation Project [37]. An overview page (http://rice.plantbio logy.msu.edu/annotation_association_analysis.shtml) provides a brief description of the procedure for identifying gene coexpression modules and contains links to pages that show trend plots for the coexpression modules for each data set analyzed. Researchers can find large-scale images of the trend plots for all modules, lists of genes from each module, and files with correlation values for all genes analyzed from each data set. Search pages allow users to query the database to explore the expression patterns of genes within a single module, within a single data set or between data sets. To enhance the functional annotation of rice genes, trend Gene expression values from a panicle and seed developmental series were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,42]. Tissues analyzed were shoot apical meristems (SAM), panicles between 0 and 3 cm long (inflorescence P1), panicles between 3 and 5 cm long (inflorescence P2), panicles between 5 and 10 cm long (inflorescence P3), panicles between 10 and 15 cm long (inflorescence P4), panicles between 15 and 20 cm long (inflorescence P5), between 22 and 30 cm long -mature pollen stage (P6), developing seed 0 to 2 days after pollination (dap; seed S1), developing seed 3 to 4 dap (seed S2), developing seed 5 to 10 dap (seed S3), developing seed 11 to 20 dap (seed S4), developing seed 21 to 29 dap (seed S5). plots for all genes covered in this study are now included on the gene annotation pages. For genes assigned to a module, the trend plot for the entire module is displayed. For genes not assigned to a module, the trend plot represents only the normalized expression values for that single gene across the treatments from the relevant experiment. In both cases, links to additional information about the module and/or parent data set are also provided.
Discussion
Gene expression data have expanded the resources available for functional annotation on a gene as well as a genomic scale. In the simplest cases, such data can help to define the tissues and conditions under which a gene is expressed. Several projects have performed correlation analyses on plant gene expression data in order to identify gene associations that may imply common functions or even regulatory relationships [6,7,8,9,10,13,17,18,19,21,22,25,26]. Many of these efforts use combined expression data sets from numerous independent experiments, and the results are typically presented in terms of complex gene association networks. In some cases, these networks are further analyzed in order to identify modules of highly correlated and connected genes.
In this study, we have performed analyses on publicly available gene expression data from a diverse collection of experiments to identify gene coexpression modules. Unlike previous studies that use combined data sets from multiple rice expression experiments [7,14,17,26], here we performed gene coexpression module analysis on expression data from individual experiments and compared it with results from a combined condition-independent data set. Our motivation in performing the condition-dependent analyses was to ensure that strong correlations apparent in select conditions were not lost when multiple diverse experiments are combined. The observation that of the genes common to both analyses, over 91% were assigned to at least one gene module from the condition-dependent analyses but only 61% were found in the condition-independent gene modules supports our reasoning ( Figure 3). Certainly, a slight change in analysis parameters could alter the numbers of genes in modules and thus shift the percentage of genes found in modules in the two analysis approaches. However, the large number of genes in many of the condition-independent modules present challenges in biological interpretation. More importantly, the common splitting of genes within a single condition-dependent module into multiple modules in the condition-independent analysis indicates that important functional associations between genes are lost through conditionindependent analysis (Table S6, Figure 3). The likely explanation for this last observation is that genes are correlated with different groups of genes within different tissues or under different physiological states. A well-defined experiment would permit the observation of one gene coexpression module, but when data from that experiment are combined with expression data from many other experiments, the correlations between the genes from that single coexpression module will be weakened and the genes in that module may be split into numerous new gene modules. Condition- The numbers listed only include those genes that passed the coefficient of variation filtering and were assigned to a module of highly correlated genes. Genes that passed the coefficient of variation filtering but that were unassigned to a module were excluded from this analysis. doi:10.1371/journal.pone.0022196.t003 independent analyses are more likely to result in gene modules with strong coexpression correlations which can obscure weaker gene coexpression relationships that occur under a subset of conditions/treatments. The obscuring effect of condition-independent expression analyses is likely to hold regardless of the algorithm or parameters used to identify gene modules. Therefore, given that our goal was to provide functional annotation to the rice gene set by identifying as many gene modules as possible, we find that the condition-dependent gene coexpression analyses are more informative.
The condition-dependent coexpression modules have been incorporated into the MSU Rice Genome Annotation Project database as an additional form of functional annotation. Of the 40,829 non-TE-related genes in the rice genome, 11,922 were assigned to at least one gene coexpression module, and 2,908 (17%) of the 17,298 rice genes that currently lack a functional description were found in at least one module. Membership in a gene module provides two distinct types of annotation to a gene. The first is association with other genes that are similarly expressed under specific conditions, and these genes may be functionally related. The second type of annotation is simply the relative pattern of expression of the gene across experimental treatments or conditions. In fact, 5,832 genes that may have been assigned to one or more coexpression modules were also found to be unassociated with any module in at least one other experiment (Table S3). The expression patterns of all genes not assigned to modules are informative as well and have been incorporated into the MSU Rice Genome Annotation Project database.
The 71 gene coexpression modules from individual experiments are diverse and will be of interest to rice researchers as these modules define sets of genes that are expressed in specific tissues or in response to various pathogen infection, abiotic stress, hormone treatments or environmental conditions (Figures 1, 2, S1, S2, S3, S4, S5, S6, S7, S8, S9, S10, S11, S12, S13). Other modules represent cultivar-specific expression differences that are apparently unrelated to experimental treatment ( Figures 2C, S1D, S4A). A statistical analysis of Pfam domain enrichment of module genes also showed that many modules have higher numbers of genes with Pfam domains related to the expected physiological state of the module, suggesting functional support for those modules (Tables 4, S4). In addition to providing annotation for genes that have been assigned to coexpression modules, the modules will be useful for formulating or supporting biological hypotheses. For example, WRKY transcription factors are often associated with regulating responses to pathogen infection [58]. A number of modules identified from biotic stress experiments contain WRKY genes, and it might be hypothesized that those transcription factors regulate the expression of other genes within those modules. Also, The green circles on the right represent the results from the conditionindependent analysis. The inner and outer circles respectively represent the genes that were assigned to modules and those that were not assigned to modules in each of the analyses, respectively. doi:10.1371/journal.pone.0022196.g003 a set of four terpene synthases and one cytochrome P450 are coexpressed in a single module from each of the Xanthomonas, Magnaporthe oryzae and S. hermonthica infection studies (Table S5), suggesting that these genes may be commonly expressed in response to a variety of biotic stresses. In contrast, numerous other chitinases, cytochrome P450s and terpene synthases were found in only one or two of these same gene modules suggesting that these genes are elicited by specific biotic stresses. When performing coexpression analysis, the choice of using a combined condition-independent data set or individual conditiondependent data depends on the goal. Additionally, the choice of parameter values will affect the numbers of modules identified and the number of genes found within those modules. The coexpression modules obtained from both condition-dependent and condition-independent data analysis are likely to be biologically relevant given that Pfam domain enrichment was observed (Tables 4, S4, S7). However, for the purposes of providing annotation to rice genes, we found that the coexpression modules identified from condition-dependent data are easier to interpret as their expression patterns are generally related to a set of treatments or tissues that are functionally related. As our goal was to provide annotation that would be intuitive to interpret, we used the normalized trend plots to guide our selection of parameters. We attempted to include as many genes as possible while obtaining gene modules with trend plots that were interpretable in a biological context. With condition-dependent analyses, we observed that genes can be assigned to multiple coexpression modules in different experiments providing numerous fine-scaled annotations that are more informative than assignment of a gene to a single module in the condition-independent method. Moreover, the multiple distinct coexpression correlations that a gene has under different physiological states can be lost or difficult to observe in condition-independent gene modules. Importantly, for an annotation project, performing gene module analysis on data from individual experiments is extensible. When new expression data become available, the results can be analyzed and added to the existing annotation. With condition-independent analysis, current coexpression results would have to be discarded and replaced with the newest analysis. Some correlations could be lost in this process, and users will find such losses to be disconcerting.
We elected to use the WGCNA method to identify coexpression modules, but the general observations from our conditiondependent versus condition-independent comparison are not expected to be different if other methods are employed. This is due in large part to the fact that most coexpression network analyses rely upon gene correlation measures, and it is the combination of expression data in a condition-independent fashion that obscures relationships that are more easily observed when condition-dependent data sets are used.
Materials and Methods
CEL files for publicly available rice expression data sets based on the Affymetrix Rice GeneChip were downloaded from either the NCBI GEO or EBI ArrayExpress [38,39] (Table S1). Arrays from individual experiments were normalized using the liwong method as implemented in the R affy package [59,60]. Quality tests were performed on the normalized array data using the Bioconductor arrayQualityMetrics package [61,62], and by examining chip trees generated by the R WGCNA package [36]. Chips that were of questionable quality were discarded. A list of all CEL files that were retained from each data set is provided in Table S1.
Probe sets from the Affymetrix Rice GeneChip were mapped to the MSU Rice Genome Annotation Project gene set (release 6.1) [37]. Individual probes were aligned to representative gene models using the vmatch alignment tool (http://www.vmatch.de). Probe sets were assigned to genes if nine or more probes from the set perfectly aligned to a single gene. Probe sets that mapped to multiple genes were discarded. If two or more probe sets mapped to a single gene, the expression value for that gene was determined by averaging the signals across the probe sets. Expression values were log 2 -transformed before being processed further. Normalized and log 2 -transformed expression values were averaged across replicate chips to generate an averaged expression value for each gene from each treatment/sample. With experiment GSE19024, biological and technical replicates were available for a subset of samples, and these were treated as simple replicates for purposes of averaging.
To reduce the number of genes for the final processing, a CV (CV = m/s) filter was applied to the averaged expression values for a single gene across a single set of conditions/treatments (condition-dependent data) or across all combined conditions/ treatments (condition-independent data) using a custom Perl script. The effect of CV filtering is to remove genes that are constitutively expressed, unexpressed or vary only modestly across experimental treatments or conditions. The CV cutoff values were determined in an ad hoc fashion with smaller CV values resulting in more genes passing the filter. Final CV values were chosen based on the number and quality of coexpression modules that were generated by WGCNA analysis ( Table 1).
The WGCNA package for R was used to identify gene coexpression modules from the normalized, log 2 transformed, CV filtered gene expression values [36]. Briefly, the WGCNA procedure calculates an unsigned expression Pearson's correlation matrix for all genes, transforms the correlation matrix by raising all values to a power ß, calculates a topological overlap matrix from the transformed correlation matrix, converts the topological overlap matrix into a dissimilarity matrix, creates a hierarchical cluster tree based on the dissimilarity matrix, and identifies gene coexpression modules from the hierarchical cluster tree using a dynamic tree cut procedure [24]. Unsigned correlations were used so that positively and negatively correlated genes could be grouped into the same cluster. The effect of transforming correlation values with the exponent ß is a form of soft thresholding that serves to strengthen strong correlation values while lessening but not discarding weak correlations. The use of soft thresholding is important for the topological overlap matrix calculation which measures the strength of two genes' correlation based on not just their direct correlation value but also the weighted correlations of all of their common neighbors [24,63]. The pickSoftThreshold function in the WGCNA package was used to determine suggested ß values. However, for most of the condition-dependent analyses, an obvious ß was not identified by this method, and in all cases, several values were tested. Higher ß values result in fewer genes with strong transformed correlation values, but with smaller ß values more genes have stronger transformed correlation values [24]. Therefore, larger ß values result in fewer genes being placed in fewer modules. Smaller ß values resulted in more genes in more modules, but with smaller the ß values, more inconsistent expression patterns of genes within individual modules were observed. The condition-independent data set used a ß value that was indicated by the WGCNA pickSoftThreshold function. A range of treecut values was also tested for module detection with larger treecut values resulting in more genes being assigned to more modules. As with the CV filter value, final ß and treecut values were chosen based on the number and quality of coexpression modules identified. All other WGCNA parameters remained at their default settings. Assessment of module quality was assisted by examining trend plots of Z-score normalized expression values for all genes in a given module (Figures 1, 2, S1 to S13). Custom Perl scripts were written to identify genes that were common to modules from both condition-independent and condition-dependent analyses.
Gene coexpression modules were tested for enrichment of genes containing Pfam domains that have been annotated within rice genes [37,64]. Statistical significance for enrichment of genes containing a particular Pfam domain was assessed using the hypergeometric distribution. A Bonferroni correction was applied to an a = 0.01 when determining statistical significance of observed Pfam domain enrichment. Figure S1 Normalized expression values of modules of genes identified from an arsenate stress study. Gene expression values from roots of rice cultivars Azucena and Bala grown in 0 ppm or 1 ppm AsO 4 were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,40]. Figure S2 Normalized expression values of modules of genes from roots and leaves in response to zeatin. Gene expression values from roots and leaves 30 and 120 min after zeatin application were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,41]. Expression data are represented here as normalized values (Z-scores). Genes responsive to zeatin treatment in roots, (A) GSE6719-blue module. Genes responsive to zeatin treatment in both roots and leaves, (B) GSE6719-brown module. Genes from leaves responsive to zeatin treatement, (C) GSE6719-green module. Genes differentially regulated in roots and leaves and also possibly regulated by zeatin, (D) GSE6719-turquoise module. Genes more strongly responsive to zeatin in roots compared to leaves, (E) GSE6719-yellow module. (EPS) Figure S3 Normalized expression values of modules of genes from seedlings in response to abiotic stresses. Gene expression values from seedlings 3 hours after stress treatments were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,42]. Expression data are represented here as normalized values (Z-scores). Genes responsive to salt stress, (A) GSE6901-blue module. Genes responsive to cold treatment, (B) GSE6901-brown module. Genes differentially regulated by drought and salt treatments, (C) GSE6901-turquoise module. (EPS) Figure S4 Normalized expression values of modules of genes identified after rice stripe virus infection. Gene expression values after infection with rice stripe virus (RSV) of rice cultivars WuYun3 and KT95 were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36]. Expression data are represented here as normalized values (Z-scores). Genes differentially expressed in WuYun3 and KT95 but not strongly regulated by RSV infection, (A) GSE11025-blue module. Genes differentially responsive to RSV infection, (B) GSE11025-brown and (C) GSE11025-turquoise modules. Genes differentially regulated by RSV infection in cultivar KT95 but not affected in cultivar WuYun3, (D) GSE11025-yellow module. (EPS) Figure S5 Normalized expression values of modules of genes expressed in gibberellin signalling mutants. Gene expression values from shoots from wild type (Taichung 65) and three gibberellin signalling mutants (gid1-3, gid2-1, slr1) were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,44]. oryzae, Xanthomonas oryzae pv. oryzicola or mock infection were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36]. Expression data are represented here as normalized values (Z-scores). Genes differentially expressed after infection with peak response after 96 hours, (A) GSE16793-blue module. Genes differentially expressed after infection with major response after 8 hours, (B) GSE16793-turquoise module. (EPS) Figure S7 Normalized expression values of modules of genes from roots and shoots after Fe and P treatments. Gene expression values from 10 day old seedlings grown with or without Fe and/or P were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,45]. Expression data are represented here as normalized values (Z-scores). Genes differentially expressed in roots in response to 2Fe and +P, (A) GSE17245-blue module. Genes differentially expressed in shoots in response to +F and +P, (B) GSE17245-brown module. Genes differentially expressed in response to the presence/absence of P, (C) GSE17245-green module. Genes differentially regulated in roots and shoots, (D) GSE17245-turquoise module. Genes differentially regulated in roots in response to Fe or P depravation, (E) GSE17245-yellow module. (EPS) Figure S8 Normalized expression values of modules of genes identified after fungal infection. Time course of gene expression values after infection with Magnaporthe oryzae strain Guy11 or mock infection were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,46]. Expression data are represented here as normalized values (Z-scores). Genes differentially expressed in response to pathogen and mock infections, (A) GSE18361-blue module. Genes differentially expressed 2 days after mock infection, (B) GSE18361-brown module. Genes differentially expressed 2 days after pathogen infection, (C) GSE18361-turquoise module. (EPS) Figure S9 Normalized expression values of modules of genes from a rice tissue survey. Gene expression values from various tissues were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,47]. Tissues sampled: germinating seed harvested 72 hour post imbibition (germinating seed); light and dark grown plumules harvested 48 h after germination (plumule 1, plumule 2); light and dark grown radicles harvested 48 h after germination (radicle 1, radicle 2); 3 day old seedling (seedling 1); trefoil stage seedling (seedling 2); less than 1 mm panicle (panicle 1); 3 to 5 mm panicle (panicle 2); 10 to 15 mm panicle (panicle 3); 40 to 50 mm panicle (panicle 4); heading panicle (panicle 5); palea/lemma 1 day before flowering (palea/lemma); stamen 1 day before flowering (stamen 1); spikelet 3 days post anthesis (spikelet); endosperm 7 days post anthesis (endosperm 1); endosperm 14 days post anthesis (endosperm 2); endosperm 21 days post anthesis (endosperm 3); shoot of seedling with three tillers (shoot); roots of seedling with three tillers (root); sheath tissues from plants with panicles less than 1 mm (sheath 1); sheath tissues from plants with panicles between 40 and 50 mm (sheath 2); leaf tissues from plants with panicles less than 1 mm (leaf 1); leaf tissues from plants with panicles between 40 and 50 mm (leaf 2); leaf tissues 5 days before heading (leaf 3); leaf tissues 14 days post anthesis (leaf 4); stem tissue 5 days before flowering (stem 1); stem tissue 14 days post anthesis (stem 2). Expression data are represented here as normalized values (Zscores). Genes expressed in shoots, mature panicles, leaf sheaths and leaf blades, (A) GSE19024-blue module. Genes expressed in spikelets and seed tissues, (B) GSE19024-brown module. Genes expressed in young and mature root tissues, (C) GSE19024-green module. Genes expressed in mature panicles and stamens, (D) GSE19024-turquoise module. Genes expressed in germinating seedling tissues, developing panicles, spikelets, shoots, roots and mature stems, (E) GSE19024-yellow module. (EPS) Figure S10 Normalized expression values of modules of genes from Rxo1 transgenic rice after bacterial infection. Gene expression values from wild type and transgenic rice containing the maize Rxo1 resistance gene after infection with Xanthomonas oryzae pv. oryzicola or mock infection were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,48]. Expression data are represented here as normalized values (Z-scores). Genes differentially expressed in wild type rice in response to X. oryzae pv. oryzicola (XOO) infection, (A) GSE19239-blue module. Genes differentially expressed in mock-infected wild type rice compared to XOO infected wild type or Rxo1 transgenic rice, (B) GSE19239-brown module. Genes responsive to XOO infection in Rxo1 transgenic rice, (C) GSE19239-green module. Genes differentially expressed in XOO infected or mock-infected wild type rice compared to Rxo1 transgenic rice, (D) GSE19239-turquoise module. Genes responsive to XOO infection in Rxo1 transgenic rice but not differentially regulated in wild type rice in response to infection, (E) GSE19239-yellow module. (EPS) Figure S11 Normalized expression values of modules of genes during aerobic germination. Time course of gene expression values during aerobic germination were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,49]. Expression data are represented here as normalized values (Z-scores). Genes with expression peaking between 1 and 3 hours after imbibition, (A) E-MEXP-1766-blue module. Genes with expression peaking after 3 hours of imbibition, (B) E-MEXP-1766-brown module. Genes differentially expressed early or late during aerobic germination, (C) E-MEXP-1766-turquoise module. (EPS) Figure S12 Normalized expression values of modules of genes during anaerobic and aerobic germination. Time course of gene expression values during anaerobic and aerobic germination were processed using Weighted Gene Coexpression Network Analysis to identify modules of highly correlated genes [36,50]. Rice seed was germinated aerobically, anaerobically, aerobically for 24 hours followed by anaerobic conditions or anaerobically for 24 hours followed by aerobic conditions. Table S8 Overlap of genes between condition-dependent gene modules and condition-independent gene modules. | 9,496.2 | 2011-07-22T00:00:00.000 | [
"Biology",
"Environmental Science",
"Computer Science"
] |
CySpanningTree : Minimal Spanning Tree computation in Cytoscape [version 1; peer review: 1 approved, 1 approved with reservations]
Simulating graph models for real world networks is made easy using software tools like Cytoscape. In this paper, we present the open-source CySpanningTree app for Cytoscape that creates a minimal/maximal spanning tree network for a given Cytoscape network. CySpanningTree provides two historical ways for calculating a spanning tree: Prim’s and Kruskal’s algorithms. Minimal spanning tree discovery in a given graph is a fundamental problem with diverse applications like spanning tree network optimization protocol, cost effective design of various kinds of networks, approximation algorithm for some NP-hard problems, cluster analysis, reducing data storage in sequencing amino acids in a protein, etc. This article demonstrates the procedure for extraction of a spanning tree from complex data sets like gene expression data and world network. The article also provides an approximate solution to the traveling salesman problem with minimum spanning tree heuristic. CySpanningTree for Cytoscape 3 is available from the Cytoscape app store.
Introduction
Graph theory is being widely used for network analysis in various fields 1 .Extraction of various kinds of subnetworks is one of the ways to identify functional modules within complex networks 2 .A tree is a subnetwork with minimal connections.Specifically in graph theory, a tree is a graph with only one path between every two nodes.In other words, any connected graph without simple cycles is a tree.Given a connected graph, which is not a tree, one can extract a tree from it by eliminating cyclic edges.A spanning tree contains all the nodes of the graph and has (N-1) edges where N is the number of nodes in the given graph.Extracting a spanning tree gets interesting when edges of the given graph have weights.In finding the minimal/maximal spanning tree, one would ideally extract the tree whose sum of weights is minimum/maximum respectively.The weight of a spanning tree is the sum of weights given to each edge of the spanning tree.There may be several minimum spanning trees of the same weight; in particular, if all the edge weights of a given graph are the same, every spanning tree of that graph is minimal.If each edge has a distinct weight then there will be only one unique minimum spanning tree.
In this paper, we present CySpanningTree, a Cytoscape 3 3 app for extracting a spanning tree from a given graph.Once the user imports a dataset, by clicking the "Create spanning tree" button of the app, a new spanning tree network is created in the network panel of Cytoscape.Historically, spanning trees are used in various applications like constructing a road network between cities with a minimum cost, as a heuristic for the traveling salesman problem (TSP), for the spanning tree network optimization protocol in networking, clustering gene expression data, etc.Three of the mentioned cases have been demonstrated in the use cases section.
Implementation
CySpanningTree is the Java implementation of Prim's 4 and Kruskal's algorithms 5 , using the Cytoscape 3 API and Java 7 for extracting a minimal spanning tree (MST).An MST for a given graph might not be unique, however for a given same Cytoscape session, the tie-breaking approach for selecting edges of equal weights is deterministic.The user gets the same spanning tree in a given Cytoscape session unless he reloads the network.This tool also has a "Create Hamiltonian cycle" button which invokes the computation of the Hamiltonian cycle 6 .For computing this cycle, it first finds an MST using Prim's algorithm and then performs a pre-order traversal on it.This pre-order traversal is a modified version of the depth-first search algorithm which results in a Hamiltonian path.Later, we connect the last node and the first node of this path to make a cycle.Users are recommended to run the Hamiltonian cycle algorithm on a fully connected graph to avoid missing of the edges while traversing.
Table 1 has the complexities of the algorithms and the uniqueness of the outputs used in the app.Prim's algorithm runs using adjacency list representation of the graph and thus implemented with a complexity O(V 2 ).Kruskal's algorithm runs using adjacency matrix of the graph and has a complexity of O(EV 2 (E+V)).The Hamiltonian cycle first calculates a spanning tree using Prim's algorithm with a complexity of O(V 2 ) and then runs depth-first search algorithm with a complexity O(E + V).
Graphical user interface
The GUI component of CySpanningTree is represented as a tabbed panel in the control panel of Cytoscape.Cytoscape takes care of loading the input network.The CySpanningTree menu (Figure 1) loads in the control panel of Cytoscape by selecting it from App menu.Currently the app runs only on connected networks.When the user tries to execute a spanning tree algorithm on an unconnected graph, an error message pops up.For weighted graphs, the user has to select the edge attribute from the drop down list (which is by default "None" that treats all edges with the same weight).
Setting the root node for Prim's spanning tree Prim's algorithm starts with a root node and hence the user is asked for the same when the Prim's Spanning Tree button is pressed.If the user enters a node that is not in the network, the user gets an error message and the program terminates.
Visualizations
The resultant MST or the Hamiltonian cycle network has the same layout as that of the input network with nodes positioned at the same location and edges scaled down.When spanning tree subnetworks are created, the corresponding spanning edges are highlighted in the input network.In Figure 2, the input network is a fully connected graph of capital cities of countries in the world, containing 203 cities and 20503 connections between them.The resultant networks: "Kruskal's Spanning Tree", "Prim's Spanning Tree" and "Hamiltonian Cycle" are connected graphs containing all the 203 cities and only 202, 202 and 203 edges respectively.Spanning trees are extracted as separate Cytoscape networks under the same network collection as shown in Figure 2.
Euclidean distance between genes g → i and g For each pair of genes, this genetic distance is calculated which gives a fully connected graph.The data set 7 has been taken from the Saccharomyces Genome Database and contains expression levels of budding yeast -S.cerevisiae with a total of 6149 genes (http:// downloads.yeastgenome.org/expression/microarray/Cho_1998_PMID_9702192/).Typically, it becomes difficult to visualize a large graph of 6149 nodes with each node connected to every other node in the graph.A spanning tree of the gene expression data makes it possible to visualize such a large network as shown in Figure 3.
• Input network: A fully connected graph of S. cerevisiae expression data • Nodes: Genes of S. cerevisiae • Edges: Euclidean distance between genes calculated using expression levels • Output network (Figure 3): Kruskal's spanning tree of the input gene expression data Although a lot of edges are removed from the network during the process of creating a spanning tree, no essential information is lost 8 .A spanning tree is a better way to visualize large networks compared to fully connected graphs.We observed that genes with similar functionalities are connected closely in the resultant spanning tree.Many clustering algorithms have been applied to gene expression data 8,9 , we are currently working on clustering using minimum spanning trees for our next release of CySpanningTree.
Use cases
In this section, we present the spanning tree results on use cases with datasets in four scenarios: gene expression matrix of gene expression data, building a cost efficient road network when all possible costs are known, an approximate solution to the travelling salesman problem and connecting a 10-home village with phone lines with minimum wiring.In each scenario, the contents of the network are introduced first and then extraction of spanning trees is demonstrated.
MST of gene expression data
The expression levels of genes when exposed to various environmental conditions are recorded at different times with different samples.This data is called gene expression data and is analyzed to extract the similarities between genes.Gene expression data ) for n genes is multi-dimensional data with each g ) for given m expression levels.Here g → i represents the i th gene and d i j represents the j th expression level of this i th gene.
This data has been simulated as a graph with nodes being genes and edges being the genetic distance between them.Genetic distance is defined as the measurement of similarity between genes.
MST on world network
This dataset 10 consists of nodes which are capital cities of all countries in the world and edges between them representing the distance in kilometers.These distances are measured using latitude and longitude coordinates of the cities (http://privatewww.essex.ac.uk/ ~ksg/data-5.html).This dataset, when imported into Cytoscape, results in a fully connected graph as the distance is calculated for each pair of capital cities.Prim's algorithm has been executed on this dataset to produce a MST network as shown in Figure 5 • Input network: Fully connected graph of capitals cities as shown in Figure 4 • Nodes: Capital cities of all countries in the world • Edges: Displacement between cities • Output minimum spanning tree: Network with minimum cost such that each city is connected.Cities separated with large distances are represented with strong edges as shown in Figure 5 Furthermore, this solution can be used for drawing a Hamiltonian cycle which is an approximation to the Travelling Salesman problem.Drawing a Hamiltonian cycle for a smaller network is discussed in the next subsection.
MST as a heuristic solution for the TSP
The TSP is a well-known combinatorial optimization problem.The goal is to find the shortest tour that visits each city in a given list exactly once and returns to the starting city.Though the problem statement looks simple, TSP is NP-complete 11 .Even though the problem is computationally difficult, a large number of heuristic solutions 12 are known due to the number of applications of this problem 13 like planning, logistics, DNA sequencing, predicting protein functions, etc.
Pre-order traversal on a minimum spanning tree is one of the heuristic solutions for TSP 5,14 .In this subsection, a Hamiltonian cycle is drawn for a spanning tree to show that the resultant cycle is a near solution to the TSP.The optimal TSP tour in Figure 9 is about 17% shorter than the Hamiltonian cycle obtained using spanning tree in Figure 8.On executing the Hamiltonian cycle algorithm on the input network, the software will create both Prim's spanning tree as well as the Hamiltonian cycle.Five nodes from the above capital city network are used for the TSP use case.
• Input network: Fully connected graph of 5 capital cities • Nodes: Capital cities of countries: USA, Brazil, South Africa, India and Italy
• Edges: Displacement between cities shown in kilometers
Connecting a 10-home village with phone lines This dataset consists of houses depicted as nodes and the edges are the means by which one house can be wired up to another.The weights of the edges dictate the distance between the houses.The task of the telephone company is to wire all houses using the least amount of telephone wiring possible.
• Input network: Houses in village depicted as graph as shown in • Nodes: Houses H 1 to H 10 • Edges: Distance between the houses • Output MST: Network which connects the houses via wires with least possible wiring.Figure 11 and Figure 12 are the spanning trees obtained using Prim's (H1 as root node) and Kruskal's algorithm, respectively.
Summary
In this paper, we present CySpanningTree app for Cytoscape 3. CySpanningTree fills an important need for many Cytoscape users and researchers in obtaining spanning trees across different types of networks.CySpanningTree makes effective use of the Cytoscape 3 API in extracting the subnetwork and creating it as a separate network.In the near future, we will be exploring MST based clustering and we are determined to explore more datasets whose spanning tree evaluation is significant.
In this research article entitled -"CySpanningTree: Minimal Spanning Tree computation in Cytoscape, the authors describe the app for Cytoscape version 3 that creates minimal/maximal spanning tree for a given network using network Prim's and Kruskal's algorithms.The CySpanningTree app appears to be useful in approximating the minimum-cost weighted perfect matching, maximum flow problems and other related issues (Supowit et al. 1980; Dahlhaus et al. 2006).The description of the proposed implementation of CySpanningTree app for Cytoscape version 3 is informative and detailed for audience.The article provides sufficient details with appropriate title and well-written abstract.
Minor Concerns Some more details on usage on practical applications are strongly suggested to include in this research article as requested by Reviewer 1 in Point 2. 1.
The definition of gene expression and generalizing gene expression data in one context is not correct in section MST of gene expression data.It is highly recommended to correct it and cite appropriate research articles defining gene expression and Gene expression data.
2.
Gene-gene interaction network reconstruction from gene expression needs to be detailed in methodology sections e.g.how edge weights are calculated and then used for calculation of Euclidean distance between genes.
3.
The usage of Genetic distance seems to be inappropriate in this context as it is a measure of the genetic divergence between species or between populations within a species.Please elaborate, if it is used in this context in research article.
4.
I would suggest making comprehensive figures for better readability e.g.(figure 1 and 5. 4. How do the authors define genetic distance?It is not clear.Is it based on correlation value of expression of genes?Please elaborate.5. Figure 5, "MST on world network" -how to use a weight; for ex., 'effective distance' between cities that is a measure of air-connectivity can be used to depict 'realistic distance' than physical distance.
6.More discussion on interpretation of figures 6,7 and figures 8,9 will be helpful to the readers.7. What is a way to verify that the solution is actually what it is 'supposed to be'?
Competing Interests: No competing interests were disclosed.
I confirm that I have read this submission and believe that I have an appropriate level of expertise to confirm that it is of an acceptable scientific standard.
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Figure 2 .
Figure 2. New networks created dynamically in Control panel.
Figure 3 .
Figure 3. Spanning tree obtained from graph of S. cerevisiae expression data; Layout: Allegro Spring-Electric layout using Allegro Layout app in Cytoscape.
Figure 4 .
Figure 4. Fully connected graph of the capital city network; Layout: Allegro Spring-Electric layout using Allegro Layout app in Cytoscape.
Figure 5 .
Figure 5. Minimum Spanning Tree of the capital city network; Layout: Allegro Spring-Electric layout using Allegro Layout app in Cytoscape.
Figure 6 .
Figure 6.Fully connected graph of 5 cities and their displacements.
Figure 7 .
Figure 7. MST of the network in Figure 6.
Figure 8 .
Figure 8. Hamiltonian cycle drawn from the spanning tree with USA as starting node.
figure 2
figure 2 may be merged into figure 1, Similarly figure 4,5,6,7 into figure 3, figure 8, 9 into figure 4 and figure 10, 11, 12 into figure 5) and brief description of figures in text as well as in legend will make help in better understanding of the examples and usage of the cySpanning trees.
Table 1 . Comparison of algorithms used in CySpanningTree.
2 + E) not unique Figure 1.User interface of CySpanningTree. | 3,635.8 | 2015-08-05T00:00:00.000 | [
"Computer Science"
] |
Spatiotemporal Evolution of Functional Structure of Urban Agglomeration in Central Yunnan
This paper selects data related to each representative industry in the Central Yunnan Urban Agglomeration from 2010-2019 as the research sample, and analyzes the functional structure of the Central Yunnan Urban Agglomeration through spatial Gini coefficient, primacy and location quotient. The research results show that: the economic development level of Kunming, a city with high primacy, is insufficient. Moreover, the development differences within the Central Yunnan Urban Agglomeration are large enough and it’s hard to drive economic recovery. The industrial agglomeration of the Central Yunnan Urban Agglomeration is not high. There are fluctuating changes in industries that depend on the natural environment, and the epidemic has a large impact on the pillar industries as well. There are overlapping industrial functions within the Central Yunnan Urban Agglomeration. The complementarity of each function is not high enough, and the degree of regional economic integration is not enough, which is not conducive to economic recovery.
Introduction
Since the 21st century, with the development of economic globalization and the increasingly close relationship between countries and regions, the role played by urban agglomerations in the national development system has become increasingly important. The development of the Central Yunnan Urban Agglomeration, as the main carrier of the economic region of Yunnan Province, has not been satisfactory in recent years. Part of the reason for this slow development is the irrational layout of industries and the lack of synergy mechanisms between cities, which has resulted in the crowding out of resources. The functional planning of the Central Yunnan Urban Agglomeration therefore requires a more precise analysis of the division of functions between the cities and an improved development strategy for the urban agglomeration, in order to address the problems in the functional structure.
Therefore, this paper proposes to study the evolution of the spatial and temporal differentiation of the functional structure of the Central Yunnan Urban Agglomeration (consisting of Kunming, Qujing, Yuxi, Chuxiong and seven cities of Honghe), using three methods: spatial Gini coefficient, primacy and location quotient. The marginal contributions of this study are: this paper will analyse the data related to each sector of the Central Yunnan Urban Agglomeration in both temporal and spatial dimensions, with the ultimate aim of promoting the optimisation of urban functions in the Central Yunnan Urban Agglomeration, and providing data support and theoretical basis for the adjustment of the functional structure of the Central Yunnan Urban Agglomeration.
Literature Review
Existing studies mostly classify cities or industries at the macro level. While at the micro level, they mostly measure the degree of the specialization of a single industry with the objective of industrial transformation and upgrading, so as to make predictions on the development trend of a single industry. For example, Ullman (1957) argued that there is a process of interaction between cities according to the division of urban functions in urban agglomerations. The cities with strong basic urban functions can export to cities with only non-basic functions, achieving coordinated development between regions [1]. On the other hand, Friedman (1963) had its definition of urban markets as well as urban system hierarchy network [2]. In terms of research methodology, Nelsen (1955) used statistical analysis and mean-standard deviation classification to classify American cities [3]. Lao, Zhang, Shen and Wang (2017) used cluster analysis to classify the functions of the middle reaches of Yangtze River Urban Agglomeration [4]. Zeng and Fang (2020) used spatial autocorrelation analysis to measure the aggregation degree of logistics industry in Guangdong Province [5]. Zeng, Li, Xing and Hu (2020) used city primacy to measure 19 urban agglomerations in China and analysed the agglomeration and diffusion effects of central cities [6]. After the system of urban agglomerations and their functional structure was gradually improved, some scholars began to analyse different levels of the functional structure of urban agglomerations in different directions. De Groot, Poot and Smit (2016) analysed specialisation, diversity, and competition through MAR externalities, Porter externalities and Jacobs externalities [7]. Ma and Zhao (2019) studied the spatial layout characteristics and industrial structure evolution of the Harbin-Changchun urban agglomeration, and the spatial correlation between industrial organization, and analyzed urban agglomeration as well [8].
Based on existing research, this paper collects and measures data on the major industries of the cities in the Central Yunnan Urban Agglomeration, comparing and analysing the data obtained, so as to make predictions on the development trends of industries and urban agglomerations, and to put forward feasible suggestions for the adjustment and development of the urban function structure of the Central Yunnan Urban Agglomeration.
Data sample selection and sources
The research object of this paper is the Central Yunnan Urban Agglomeration. The data used for each city are sourced from the Yunnan Statistical Yearbook from 2010-2019, mainly including the number of people employed in eight industries in cities from both the Central Yunnan Urban Agglomeration and Yunnan Province, the total employment in cities from both the Central Yunnan Urban Agglomeration and Yunnan Province, the population in cities from both the Central Yunnan Urban Agglomeration and Yunnan Province, and GDP in cities from both the Central Yunnan Urban Agglomeration and Yunnan Province.
Model Specification
In order to better reflect the characteristics of the functional structure of the Central Yunnan Urban Agglomeration and the trend of spatial and temporal divergence, three models are used in this study to build the research process: Firstly, the spatial Gini coefficient formula is used to measure the spatial Gini coefficients of the major industries in the Central Yunnan Urban Agglomeration and the five cities/autonomous prefectures it contains, and the data is visualised in a graph. The spatial Gini coefficient is calculated as follows: Where: is the spatial Gini coefficient of the industry, is the proportion of employment in a particular industry (i.e. flowers, seedlings, construction, etc.) in the region i to employment in that industry in the province, and is the proportion of employment in that region to total employment in the province.
Secondly, this paper uses economic primacy and city primacy to reflect the economic scale structure, city scale structure and population concentration of the Central Yunnan Urban Agglomeration. It is generally believed that a city's degree of primacy is less than 2, indicating a normal structure and appropriate concentration; the degree is greater than 2, indecating there is a tendency of structural imbalance and over-concentration. The calculation formula is as follows: Where: S t is the economic primacy of city , G is the economic scale (GDP) in Kunming, and G t is the economic scale of city .
S i is the city primacy of city i, P is the population scale of Kunming, the first city in the Central Yunnan Urban Agglomeration, P i is the population size of city i.
Finally, this study applies the location quotient to measure the degree of development, specialisation and the economic support base that each city in the Central Yunnan Urban Agglomeration can bring to the region by industry. It is calculated as follows: Where: Lq ij is the location quotient in industry in city i; G ij is the number of workers in industry in city i; G i is the number of workers in city i; G j is the number of workers in industry in Yunnan Province; G is the number of workers in Yunnan Province.
Based on the spatial Gini coefficient
Using Model (1), the spatial Gini coefficients for each industry in the Central Yunnan Urban Agglomeration and the five cities/autonomous prefectures can be obtained for the decade 2010-2019.
On the whole, the spatial Gini coefficients of all industies in the urban agglomeration of central Yunnan Province in China are mostly on a stable trend. Meanwhile, the spatial Gini coefficients of the health, social security, wholesale and retail, accommodation and catering, manufacturing and construction industries all remain at a low value year-round.
At the same time, it can be found from Fig.1 that the main industries with high spatial Gini coefficients among the various industries in the Central Yunnan Urban Agglomeration are mining and agriculture, forestry and fishery. Overall, the spatial Gini coefficient of the mining industry is perennially at the highest value among all industries and is the main industry for the development of the Central Yunnan Urban Agglomeration. But its value fluctuates widely and could be closely related to the discovery or depletion of mineral resources. Secondly, the spatial Gini coefficient of agriculture, forestry, animal husbandry and fishery is perennially in the second position in the Central Yunnan Urban Agglomeration, but its value declined steeply in 2017-2019, which may be related to the significant decrease of resources for agriculture, forestry, animal husbandry and fishery in the Central Yunnan Urban Agglomeration. On the other hand, the spatial Gini coefficient of the electricity, gas and water production and supply industry in the Central Yunnan Urban Agglomeration has been in a low state from 2010-2019, but its value suddenly increased in 2019. The situation may be due to the remarkable process Kunming has achieved in advancing this aspect of the electricity pillar industry in Yunnan Province, mainly hydropower, which has led to a large demand for employment in this industry, and therefore a sharp increase in employment driving the spatial Gini coefficient higher [9]. The analysis above shows that in the past decade, mining and agriculture, forestry, animal husbandry and fishery industries have played a pivotal role in the development of the Central Yunnan Urban Agglomeration, relying on a high degree of spatial aggregation and a relatively mature industrial system to achieve the efficient deployment of production factors and the rational use of resources. However, the mining and agriculture, forestry and fishery industries are both extremely dependent on natural resources with non-renewable characteristics, indicating that there will be significant bottlenecks in these two industries. According to the calculation results, the degree of concentration of these two industries in the Central Yunnan Urban Agglomeration has shown a decreasing trend from 2016 to 2019. It indicates that some capital and labour have already withdrawn from these two industries, which greatly reduces the long-term dependence of the Central Yunnan Urban Agglomeration on these two industries. This greatly reduces the possibility of the Central Yunnan Urban Agglomeration relying on these two industries for long-term development.
Spatial Gini Coefficient Trends by Sector in Central Yunnan Urban Agglomeration
Therefore, there is a need to adjust the current state of the industrial structure of the Central Yunnan Urban Agglomeration. Firstly, the development patterns of mining and agriculture, forestry and fishery need to be adjusted. As once leading industries, the two are important components of the industrial structure of the Central Yunnan Urban Agglomeration. So rational planning of resources, reducing the proportion of the two in the industrial structure and providing some appropriate protection policies for the relevant resources are planed not only for the long-term development of the industries, but also in helping to buffer the impact of industrial restructuring. Secondly, there is a necessity to find new industries that can be developed in the long term. In terms of data, the concentration of the electricity, gas and water production and supply industry in the Central Yunnan Urban Agglomeration has increased significantly, while there is a decline in the concentration of the mining and agriculture, forestry and fishery industries, from which I could interpret as a subtle shift in the market. In essence, these three industries are all dependent on natural resources, but the biggest difference between the electricity, gas and water production and supply industry and the first two is that the resources they rely on are highly renewable, which means that they have the potential to develop as dominant industries in the long term. At the same time, the fact that it has the same origin as the first two industries can help it to absorb the capital and labour lost by the first two industries, which is a very important advantage over other industries. It is also an important opportunity for the Central Yunnan Urban Agglomeration to adjust its industrial structure for sustainable development on its own.
Economic primacy
By substituting the relevant data into model (2), Table 1 was obtained. According to Table 1, the economic primacy value of the Central Yunnan Urban Agglomeration, represented by Kunming City, increased slowly from 2010 to 2019, and reached a maximum value of 1.125, which is less than 2 in 2019. It indicates that the economic development structure of the Central Yunnan Urban Agglomeration is normal and the degree of industrial agglomeration is not high. The highest value is less than 1.5, which indicates that Kunming's economic development is average. But as the first city, Kunming has weak economic impetus to the Central Yunnan Urban Agglomeration and may not be able to drive the development of the whole region. From the analysis ablove, it can be recognized that the scale distribution of the Central Yunnan Urban Agglomeration belongs to the moderate first place distribution of the first-place distribution type. The city scale hierarchy is unbalanced, and the first place of the central city needs to improve the city scale structure. Break the single mode of expanding the city population, and instead attract advantageous industries and high-quality talents to further enhance and strengthen regional agglomeration and lead the functional development of the surrounding cities.
City primacy
By substituting the relevant data into model (3), Table 2 was obtained.
In recent years, due to the export constraints, national economic strategies have gradually tended towards endogenous growth. Based on this situation, the economic drive of the country's top city on the surrounding urban agglomerations is particularly important. According to Table 2, the city scale primacy of cities in Central Yunnan represented by Kunming fluctuates from 2010 to 2019, but the overall average is greater than 2. This indicates that the cities in Central Yunnan represented by Kunming are the gathering place of talents and have higher economic vitality. But it also shows that the development of population scale in Central Yunnan is uneven, with large differences in local development. At the same time, the spatial structure of the Central Yunnan Urban Agglomeration has been profoundly altered by the interaction of various factors. In the Central Yunnan Urban Agglomeration, Kunming, the first place of the central city, has fluctuated over the decade, but it is still at its highest level. As the large-scale development of Kunming, the cost from the surrounding provinces and regions is far less in terms of resources than in terms of economic-driven benefits. It slows down the predatory effect of the first city as well. Therefore, the development of the first city is still beneficial to the economic growth of the nearby provinces and the gradual formation of a high spatial-geographical concentration of the economies of the various provinces will lead to an increase in the degree of openness to the outside world, which will help to promote the development of the Central Yunnan Urban Agglomeration as a whole.
Analysis of the urban location quotient of the Central Yunnan Urban Agglomeration from an industry perspective
By substituting the relevant data into model (4), Table 3 was obtained.
On the whole, Kunming has a relatively balanced overall development, except for heavy industry, and exceeds the other sub-centres; Qujing and Yuxi have a relatively stronger development; and the overall level of Chuxiong Prefecture needs to be improved. The city with the highest level of specialisation and agglomeration in the construction industry is Kunming, which remains at a high level overall. As the Wujiaba area in Kunming has undergone large-scale renovation and construction of a large urban central park in recent years, it attracts a large number of real estate industries [10]. At the the same time, the real estate and construction industries are cooperative parties, so the renovation of this area has had a positive impact on the construction industry in Kunming in some ways. highest and lowest levels of specialisation are in Chuxiong and Kunming in comparison. In terms of agriculture, although Yunnan is a traditional agricultural province, the agriculture, forestry and fishery industry is not prime and dominant in the Central Yunnan Urban Agglomeration. The main reason is that some places of Yunnan's specialities, such as Lincang for pu-erh tea, Yuanmou County for high quality vegetables and Dehong for small grain coffee beans, are not the part of the Central Yunnan Urban Agglomeration. The Central Yunnan Urban Agglomeration is not dominated by modern or traditional agriculture.
Analysis of the complementarity of functional structure sectors from the perspective of urban agglomerations
Although there are minor changes in the mainstream industries of the cities, but in general, the mainstream industries of the cities do not vary significantly and the functional structure of the Central Yunnan Urban Agglomeration is relatively stable. We have analysed the mainstream industries in each city with a frequency greater than 50% in each year as the mainstream industries in the Central Yunnan Urban Agglomeration over the past 10 years, and plotted the following table： Table 4. Analysis of mainstream industries in Central Yunnan Urban Agglomeration.
Kunming P P P P Qujing P P P P Yuxi P P Chuxiong P P P The highest industry complementarities are in the health and social security, accommodation and catering. In terms of the health and social security, Chuxiong can export industries to Kunming, Qujing and Yuxi. Kunming, which is the most prosperous and commercial in the city cluster, has the highest concentration of accommodation and catering, which can be exported to other sub-centres, thus sharing resources.
Industries with relatively high complementarity are the wholesale and retail trade, mining, construction, and electrical, gas and water production and supply. Kunming is the central city of the Central Yunnan Urban Agglomeration, with convenient transport links for the flow of goods. In addition, Kunming and Yuxi are adjacent to each other, and it allows for the formation of a larger "Kunming-Yuxi" wholesale and retail market centre, which can better radiate the Central Yunnan Urban Agglomeration and export to surrounding cities. The mining industry and the electrical, gas and water production and supply industry are relatively highly complementary. But there is a slight overlap of urban functions in the mining industry and the electrical, gas and water production and supply industry in Qujing due to the over-concentration and high specialisation of the mining industry in Qujing. Additionally, Qujing is rich in water, gas and electricity. However, Chuxiong Prefecture is also in a position to have an advantage in similar industries within the Central Yunnan Urban Agglomeration. Among the four cities, there are three cities, Kunming, Qujing and Yuxi, having the manufacturing as the mainstream urban industry, with a large overlap and low industrial complementarity. Agriculture, forestry and fishery industry, due to the natural environment, are basically not self-sufficient in the Central Yunnan Urban Agglomeration and need to be imported from neighbouring cities in Yunnan Province.
Optimising the functional structure of the Central Yunnan Urban Agglomeration
According to the data above, although the wholesale and retail and accommodation and catering are both mainstream industries in Kunming, neither industry has exceeded a location quotient of 1.5 at the same time in the past ten years, becoming an industry with relatively high comparative advantages at the same time. Furthermore, on the whole, there is still room for development in wholesale and retail trade compared to accommodation and catering. The healthy development of the tertiary industry represented by the wholesale and retail industry and the accommodation and catering industry can balance the over-tourism of the tertiary industry in the Central Yunnan Urban Agglomeration and comprehensively coordinate the economic structure [11]. Further, the enhancement of the restaurant industry will have a positive impact on the service trade [12]. When the location quotient of both can exceed 1.5 at the same time, it not only allows the wholesale and retail industry and the accommodation and catering industry to be fully developed, but also increases the core competitiveness of Kunming as the central city of the Central Yunnan Urban Agglomeration. Increase the diversification of business models, brand building characteristics, and increase the proportion of trade. Use Kunming as a window to seek a larger international market.
According to Table 4, Kunming and Qujing each have four mainstream industries, making them the two cities with the most mainstream industries in the Central Yunnan Urban Agglomeration. However, there is an overlap in the division of urban functions in manufacturing and construction. On the other hand, the construction and manufacturing industries in the Central Yunnan Urban Agglomeration are also largely driven by Kunming and Qujing.
For Kunming City, it can turn to the investment in manufacturing industry and increase the investment in high-tech industry to alleviate the high overlap rate of manufacturing industry in Central Yunnan Urban Agglomeration. Additionally, the other target is to solve the problem of resources for reasonable coordination, increase the core competitiveness in the central city, improve the economic status and increase the economic radiation to Central Yunnan Urban Agglomeration.
At the same time, as Qujing City and Chuxiong Prefecture in the mining industry and electrical, gas and water production and supply industry have a certain degree of overlap, it's also necessary to reduce investment and constructio of Chuxiong Prefecture in these two industries.
Integrated and comprehensive planning of the functions and positioning of the Central Yunnan Urban Agglomeration
The optimisation of the functional structure of the cities in Central Yunnan should be planned in a comprehensive manner from the height of the province. Considering regional strengths, urban development needs, development conditions and development bases, comprehensively integrating economic, social, resource and environmental factors, the scientific and reasonable adjustments should be formulated to the industrial countermeasures of the Central Yunnan Urban Agglomeration to guide the relative transformation and development of its functional structure. Position the urban function of cities to make it coordinate and sustainable, and establish a corresponding city image. Ensure that the functions of the cities are subordinated to the general situation of healthy urban development in the province, which is conducive to the sustainable development of the regional urban economy and the improvement of people's living standards.
Ensuring the government's macro-regulatory role over the functional structure of individual cities
The government's macro-control should be strengthened. Firstly, consolidate measures for the stability of advantageous industries. Secondly, support the development of some industries in light of the current situation, so as to reduce the impact of the new epidemic on the social economy and even create conditions for a good multi-polar development in the future. According to the analysis based on the spatial Gini coefficient and location quotient, the relevant protection policies should be introduced for the mining industry, agriculture, forestry and fishery industriy.
Use of market regulation function to promote industrial transfer
Market regulation is an important means of rational allocation of resources, and industrial transfer is one of the most direct and effective means of restructuring the functions of cities. Giving full play to the regulating function of the market and promoting the reasonable transfer of the main industries in central Yunnan cities within the national region will help the industries in central Yunnan cities to construct a scientific industrial system and optimize the functional structure of the cities. As domestic production costs, especially labour costs, are increasing, the market, guided by the law of value, spontaneously regulates the operation of the economy. Labour and some capital-intensive industries are gradually shifting from the eastern coastal areas, where production costs are high, to the inland areas, where labour resources are abundant and costs are low, whch is an inevitable trend. The transfer of industries between cities in China is more active. By making full use of the regulating function of the market to actively guide the transfer of industries within and between cities under macro-control, the system of city functions will be more adapted to the requirements of healthy urban development.
Improve competitiveness by highlighting the advantageous functions of the first city
Under the background of the economic globalisation, the development of cities into a global circumstance depends not on their own shortcomings, but on the dominant conditions. Highlighting the city's advantageous functions will enable it to take its place in the global development for the big data. It is particularly important for Kunming, the first city in the Central Yunnan City Agglomeration, to highlight its own city's superior functions. From the perspective of primacy analysis, it shows that Kunming has an average level of economic development and lacks the ability to drive the development of the Central Yunnan City Agglomeration as the first city. Therefore, the Central Yunnan Urban Agglomeration needs to improve the economic primacy of Kunming as the first city in central Yunnan, and combine the industrial structure of Kunming with relevant national policies to meet the needs of Kunming's green industrial development. Relying on Kunming's resource advantages, vigorously develop Kunming's advantageous and characteristic industries, improve the economic vitality of the Central Yunnan Urban Agglomeration, and promote Kunming to achieve the "win-win" goal of ecological construction and economic development as soon as possible. It will also enable Kunming to realise the 13th Five-Year Plan as soon as possible and become an environmental and ecological home for harmonious development. | 5,863.2 | 2021-01-01T00:00:00.000 | [
"Economics"
] |
Towards Axion Monodromy Inflation with Warped KK-Modes
We present a particularly simple model of axion monodromy: Our axion is the lowest-lying KK-mode of the RR-2-form-potential C2 in the standard Klebanov-Strassler throat. One can think of this inflaton candidate as being defined by the integral of C2 over the S 2 cycle of the throat. It obtains an exponentially small mass from the IR-region in which the S2 shrinks to zero size both with respect to the Planck scale and the mass scale of local modes of the throat. Crucially, the S2 cycle has to be shared between two throats, such that the second locus where the S2 shrinks is also in a warped region. Well-known problems like the potentially dangerous back-reaction of brane/antibrane pairs and explicit supersymmetry breaking are not present in our scenario. However, the inflaton back-reaction starts to deform the geometry strongly once the field excursion approaches the Planck scale. We derive the system of differential equations required to treat this effect quantitatively. Numerical work is required to decide whether back-reaction makes the model suitable for realistic inflation. While we have to leave this crucial issue to future studies, we find it interesting that such a simple and explicit stringy monodromy model allows an originally sub-Planckian axion to go through many periods with full quantitative control before back-reaction becomes strong. Also, the mere existence of our ultra-light throat mode (with double exponentially suppressed mass) is noteworthy. December 14, 2015 ar X iv :1 51 2. 04 46 3v 3 [ he pth ] 6 D ec 2 01 6
In [38] it was shown that models of axion monodromy inflation in the complex structure moduli sector of Calabi-Yau 3-and 4-folds require a significant level of tuning to avoid excessive backreaction and the destabilization of Kähler moduli. The required level of tuning can only be achieved in 4-folds which further complicates the model. These difficulties can be avoided if Kähler moduli are stabilized using non-geometric fluxes [39][40][41]. However, it remains a challenge to implement a consistent hierarchy of scales in the resulting models.
Given the technical difficulties encountered in most constructions of axion monodromy inflation, it would be desirable to realize as minimal a model of axion monodromy inflation as possible. In such a simple construction one may hope that questions regarding the consistency and detailed phenomenology can be addressed explicitly and quantitatively. This is what we set out to do in this work. Here, we present a simple model of axion monodromy which is based on the standard Klebanov-Strassler-throat [42] (i.e. the deformed conifold) with shrinking S 2 .
Our axion is the RR-2-form C 2 wrapped on the homologically trivial S 2 , similarly to some of the settings in [31]. We do not need to include branes in our setup, the main point being that the axion acquires its monodromic potential from the homological triviality of the S 2 (in contrast to models where the potential is due to the tension of the NS5-brane). Thus we do not need to include anti-branes either and therefore evade the dangerous brane/antibrane back-reaction described in [34,35]. We note that our results might also be useful in the context of recently proposed Relaxion-models [43][44][45][46][47][48][49][50][51][52].
We find that the mass of the lightest 4d-Kaluza-Klein mode is lighter than the next heavier mode by a relative warp-factor which makes it an interesting candidate for single field inflation. Thus the inflaton potential is suppressed by warping [53] without the need for an additional tuning. Since this is due to the S 2 ending in the infrared-region we need a second throat into which the S 2 can bend around in the UV such that its second end lies in an infrared region as well. Such a geometry has been constructed in [54,55] which we very briefly review in Section 2.
This paper is organized as follows: In Section 3 we calculate the IR-localized 5d-massterm, finding that Λ ∼ 1/R where R is the typical radius of the KS-region. In Section 4, starting from the 5d-effective model we perform a Kaluza-Klein-reduction along the radial coordinate of the throat, thereby obtaining the effective 4d-theory with an infinite tower of KK-modes with the above mentioned mass-suppression of the lightest mode. In Section 5 we compare the energy-densities of the inflaton with those stabilizing the throat, concluding that an explicit numerical back-reaction study is necessary to make statements about the stability of the KS-throat at large field excursion 2 . In Section 6 the parametrization of the fully back-reacted inflaton mode in the KS-throat is given while the differential equations that need to be solved are listed in Appendix A. We draw our conclusions in Section 7.
2 Note that this is in contrast to a more optimistic claim of an earlier version of this paper.
The Double Throat
Let us briefly review the construction of the double throat (see Figure 1) following the discussion in [55] 3 . The conifold can be described as the subset of C 4 solving The conifold singularity sits at x = y = u = v = 0. We can construct a two-conifold-setup by replacing x with a polynomial W (x) in the conifold equation (1). We take W to have two simple roots at x ∈ {a 1 , a 2 }: If g = 0 this gives a curve of A 1 -singularities parametrized by x. Blowing up the singularity gives a curve of P 1 's. Setting g = 0 there is still a family of S 2 's related in homology. After a geometric transition [54] the system is deformed by means of a polynomial f 1 of degree one, to give two deformed conifolds with shrinking S 2 : This is precisely the geometry we will use. Figure 2: The geometry close to the tip of the throat.
A Simple Geometric Setup and Reduction to 5d
Consider the standard KS-throat with a blown up S 3 KS but trivial S 2 -cycle. Due to the homological triviality of the S 2 there is no harmonic 2-form and thus no massless axion c = S 2 C 2 . Our axion will hence be the (massive) lightest KK-mode of C 2 .
As a first approximation, let us take the geometry of the compact space to be simply of constant radius R, which is closed by one half of a three-sphere (≡ S 3 1/2 ) in the IR. In the UV the S 2 bends around into a second throat such that it is closed in the IR on both sides as depicted in Figure 1. This is crucial since we would otherwise generate a UV-mass-term.
Let y be the radial coordinate such that y = 0 at the boundary of S 3 1/2 (see Figure 2). Our starting point is the type IIB supergravity action in Einstein frame (see [56], ch. 12.1): where F 3 = dC 2 and H 3 = dB 2 are the three-form field-strengths and we have restricted ourselves to constant dilaton e φ ≡ g s and vanishing C 0 . We now expand: where we take ω 2 to be the canonical volume-form of S 2 (normalized to ω 2 = (V ol S 2 ) −1 * 2 1).
We now want to derive the effective 5d-action which we will then treat as an effectively 5d Randall-Sundrum-model [57,58] in Section 4 (see [59] for the 5d-description of the throat).
First we derive the bulk-term. Thus we plug the above into the 10d-(Einstein-frame)action S IIB ⊃ − gs 4κ 2 10 dC 2 ∧ * dC 2 and get a bulk kinetic term − g s 4κ 2 10 M 5 Next let us calculate the contribution of the boundary S 3 1/2 . Since S 2 is trivial (e.g. at y = 0, with F 3 = dC 2 . Neglecting the warping, the lowest energy configuration is where the fieldstrength F 3 is equally distributed over S 3 1/2 . Hence we make an ansatz where ω 3 is the canonical volume form of the three-sphere (i.e.
Plugging this into the 10d-action we get a boundary mass term where we have again neglected the effect of warping on S 3 1/2 . Going over to a canonically normalized 5d-field ( gs we get a 5d-action Therefore the localized mass-term is essentially Λ ∼ R −1 where R is the typical transverse size of the throat which in this case coincides with the length-scale over which the throat contracts.
KK-Reduction on the Effective 5d-Throat and the 4d Action
The 5d action derived in the previous section can now be reduced to an effective 4d action containing an infinite tower of 4d-KK-modes. We now treat the throat as an effectively 5-dimensional Randall-Sundrum-model [57][58][59].
We now let y take values on a strip of length L choosing orbifold identification y ∼ = −y and y ∼ = y + 2L (note that in this case we need to double Λ, that is Λ = 8/πR, in order to get the physical boundary-condition for the 5d-field).
Note that the delta-potentials come from enforcing the appropriate boundary conditions on χ (not on f ). The general solution (a special case of the more general situation considered in [61]) now takes the form where J n and Y n are the Bessel functions of first and second kind respectively.
From the form of the potential we immediately deduce the existence of a single (UV-) bound state and wave solutions of higher energy (mass) that are exponentially suppressed in the UV. Note that the bound state solution can be determined exactly in the case where Λ = 0: which simplifies to χ = const. after imposing boundary conditions. This is of course the constant mode of zero mass which can be immediately read of from (14).
The mass-condition follows from the two boundary conditions (∂ y χ(0) = Λ 2 χ(0) and ∂ y χ(e −kL ) = 0) and reads We will now focus on the case r c ≡ 1 k L (which is the interesting case of strong warping).
For the bound-state solution we expect a small mass (m k) for which we can use the small argument approximations of the Bessel-functions to arrive at Remarkably this mass is exponentially suppressed by the warp factor (thereby a posteriori justifying our small argument approximation). It is crucial to realize that this is not the usual hierarchy induced by warping in Randall-Sundrum models [57] but is rather a suppression 'on top of that' since our metric conventions are such that g IR µν ≡ g µν (y = 0) = η µν . The zero-mode profile takes the following form: The higher KK-modes (with 1 m k e kL ) are obtained by noting that
such that the mass condition is approximately
The solutions interpolate between the zeros of the two Bessel-functions (j 1,n and j 2,n ), that and asymptotically (that is m n k, Λ) which are the usual KK-masses but with L replaced by the curvature radius r c ≡ k −1 . The bound-state and the first excited states are plotted in Figure 3.
Using the 4d-Planck mass [57] one immediately sees the double exponential suppression of the bound-mode: Note that this agrees with the expression for the axion-potential in equation (4.76) of [18] where the potential comes from the NS5-DBI-action.
This behavior could have already been anticipated from the form of the potential (16): The bound-state-solution approaches a constant in the UV while the positive delta-potential in the IR leads to a dip in the IR. It therefore gets its mass from the IR while its kinetic term lives in the whole bulk (concerning the kinetic term arguments along these lines have already been given in [18], Sec. 4.3.2). This leads to the already mentioned 'double'-suppression.
The higher KK-modes are the solutions to Schrödinger's equation (15) that oscillate in the IR-region 0 < y r c and fall off exponentially towards the UV due to the ∼ 1/z 2 -term in the potential (16). This leads to the modified KK-mass-formula (25).
Note furthermore that m 0 (more precisely its upper bound) is not particularly sensitive to the value of Λ: where Let us pause here and highlight what we have found: The lightest KK-mode of the RR-2-form C 2 on the KS-throat with trivial S 2 -cycle is exponentially lighter than the next higher mode in the case of strong warping. This makes it an ideal candidate for single-field chaotic inflation since we can safely ignore the higher modes.
Energy Density at the boundary of S 3 1/2
It is important to check that the energy density at y = 0 on the cylinder is the same (at least up to O(1)-factors) as the one on S 3 1/2 . On the cylinder we have which implies that ε cyl ε S 3 Therefore the energy-densities are exactly the same. Note that this were also true if we had chosen any other eigen-mode of the 5d-Laplacian since the identity Λ = ∂ y χ n (0)/χ(0) is simply the boundary condition for the y-profile of any mode χ n .
Inflaton Energy Density vs F 3 -Flux Energy Density
Since we have a model of single-field large field inflation we have to make sure that the field excursion of the inflaton does not back-react in a way that destabilizes the throat.
The flux energy density can be calculated from the type IIB Supergravity action where ω 3 is the appropriately normalized volume form on S 3 KS and M is the F 3 -flux on S 3 KS stabilizing the throat. Ellipsis indicate terms that integrate to zero over S 3 KS .
where R is the radius of the S 3 KS (which we identify with the S 3 1/2 radius). In the second step we have used that The inflaton energy density (using equations (10), (11) and the explicit form of the bound mode (22)) is given by where α measures the 4d field excursion in 4d-Planck units (equivalently the 5d excursion in 5d-Planck units). The ratio of the densities therefore satisfies Therefore in the interesting regime of large field, α 2 1, the back-reaction on the ambient geometry cannot be neglected. The full non-linear equations of type IIB Supergravity have to be considered to quantify this back-reaction.
The Ultra-light Mode in the KS background
In the following we would like to describe the ultra-light mode in the full Klebanov-Strassler (KS) geometry in order to address questions of back-reaction. Since back-reaction effects take place at the tip of the throat only, where the metric is known, this can be done explicitly.
To this end we will specify an explicit ansatz that describes our mode in the KS-background and derive the equations of motion. Obtaining the full solution is an involved numerical task, that will be left for future research.
A Simple Prescription for Obtaining the Back-reacted Potential
Before turning to the relevant equations of motion, let us discuss how the effective backreacted potential in 4d can be obtained without having to solve complicated time-dependent equations of motion. As we will see, the effective 4d potential can be efficiently extracted by considering static and homogeneous field profiles φ = φ(τ ).
Let us parameterize the effective 5d action as follows, where we use the for now arbitrary radial coordinate τ which does not necessarily measure physical distances (i.e. g τ τ need not be unity). Here the function X(τ ) parameterizes the varying volume of T 1,1 and its 2-cycle which appear in the dimensional reduction from 10d to 5d and e 2A is the warp factor. We have not written out explicitly any terms beyond quadratic order in φ. These are included in L int which we assume to take significant values only near the IR. Clearly there is no static homogeneous solution to the equations of motion with the boundary conditions φ(0) = 0 = ∂ τ φ(τ U V ) other than the trivial solution φ = 0. This is expected as we know that the lowest lying mode obtains a non-vanishing potential from the 4d-perspective and can hence not be static. However, if a source j is inserted at the UV boundary, a non-trivial solution is obtained as the UV-boundary conditions are altered to For given source j there is hence a non-trivial static profile φ(τ ) that solves the (non-linear) equations of motion. Intuitively the source j sets the field excursion by applying a restoring force against the potential slope. Let us parameterize the field-excursion by the value φ U V ≡ φ(τ U V ). Then to each value of the source j there is an associated field excursion φ U V and (on-shell) we can hence interpret the source j as a function of the field excursion, It should be noted that in order to obtain this function j(φ U V ) explicitly, the non-linear equations of motion have to be solved numerically. The function j(φ U V ) is then a complicated non-linear function that is known only numerically.
Let us now change perspective and analyze the same problem from the effective 4d point of view. The 4d action is with axion decay constant f and potential V (φ U V ). At this stage, the potential V (φ U V ) is unknown. Again, there is a static configuration at field excursion φ U V if Because both the 5d point of view as well as the effective 4d point of view should give the same answer, the potential V (φ U V ) can be inferred by comparing (39) with (42). Finally, we have obtained the desired simple prescription to read off the effective 4d-potential from a static numerical solution of the non-linear bulk equations of motion with boundary conditions The crucial advantage is that there is no need for an explicit dimensional reduction of the higher-dimensional action to 4d.
Let us now specify to the case of strong warping, approximately constant field-profile φ(τ ) and g τ τ ≈ const ∼ k −2 at large τ , where k is the inverse curvature radius of the effective 5d geometry. Then the kinetic term is dominated by the UV-region and it follows that where we have dropped overall factors of O(1) that are not affected by back-reaction. In the case of the KS-throat one has that k −2 ∼ g s M α and we call m 2 wKK ≡ (g s M α ) −1 the warped KK-scale which is the mass-scale of KK-modes that are localized at the tip of the throat.
Then it follows that the effective potential in 4d can be expressed as which corresponds to a source j that is linear in the field excursion φ U V .
The Type IIB Equations of Motion
Having learned how to extract the effective potential from a solution to the equations of motion we now derive the explicit equations of motion that need to be solved eventually. We start with the String frame equations of motion and Bianchi-identities of type IIB Supergravity (for now omitting the Einstein equations): In practice we will work with F 1 = dC 0 , H 3 = dB 2 andF 3 = F 3 − C 0 ∧ H 3 and specify an ansatz for F 3 , B 2 and C 0 such that dF 3 = 0. We will not work with a four-form potential and specify an ansatz directly for F 5 . Furthermore we redefine the dilaton Φ as e Φ ≡ g s e φ and define g s to be the value that e Φ approaches in the UV. and v 2 ≡ (g 3 , g 4 ) T while it has trivial action on g 5 [63]. As a result, the following 2-forms and symmetric 2-tensors are invariant under the symmetries of the deformed conifold 7 :
The Ansatz and Boundary Conditions
Here, g i g j ≡ 1 2 (g i ⊗ g j + g j ⊗ g i ). Including the radial direction parameterized by the coordinate τ , one may further allow for a term proportional to dτ g 5 in the 6d metric since g 5 is invariant under the symmetries of the deformed conifold. This leads to the 6d metric re-parametrization ψ −→ ψ + λ(τ ), one has that g 5 −→ g 5 + λ (τ )dτ . Under such a reparametrization the 6d metric is not invariant but again takes the form of (53) with different coefficients. In particular it can be checked that A non-vanishing function d(τ ) can thus be gauged away by a suitable re-parametrization and we fix the gauge by setting d(τ ) ≡ 0.
Consequently we choose the ansatz for the 10d metric. The zehnbein one-forms are with radially varying functions A, B, C, D, E 8 . One can check that this choice of zehnbein one-forms reproduces all the terms in (53) in a sufficiently general way.
We generalize the KS-ansatz to F 3 = M α 2 g 5 ∧ g 3 ∧ g 4 + d(δC 2 ) and H 3 = dB 2 with It then follows by virtue of equation (48) that Furthermore we allow for radial profiles of the axio-dilaton in a convenient parametrization The IR boundary conditions are for the axio-dilaton and p-form fields, while we choose to parametrize the metric function B by B(τ ) ≡ tanh(τ )B(τ ) and set Because the harmonic two-form of T 1,1 is proportional to g 1 ∧ g 2 + g 3 ∧ g 4 we define the field excursion ψ by imposing the UV boundary condition Furthermore all other functions are required to take their KS-values. The set of IR boundary conditions (60) follows from demanding that field strengths be finite at the tip and (61) is a consequence of demanding the tip-topology to still be an S 3 . The axio-dilaton is stabilized in the UV by ISD-fluxes on other cycles that are not relevant to the local KS-throat and we implement this by demanding that Finally let us note that we have parametrized our ansatz such that at ψ = O(1) backreaction effects become strong. The C 2 -axion φ U V ≡ 1 2πα S 2 δC 2 (τ U V ) with natural periodicity φ U V −→ φ U V + 2π can be expressed in terms of the field ψ by φ U V = 2M ψ and the axion-decay constant 9 f can be estimated to be where ω 2 is the harmonic 2-form of T 1,1 normalized to S 2 ω 2 = 1 and B and C are dimensionless radial functions that appear in the KS metric (see (56)). Here we have assumed that the 2-cycle is of constant size as one passes from one throat to the other through the compact CY. The canonical field excursion in 4d can then be related to the variable ψ by To get an upper bound on the throat length that is needed to generated the desired hierarchy let us assume that warping is the only effect that lowers the mass scale of the ultra-light mode 10 . Then in order to achieve Hence, back-reaction becomes strong at field excursion φ c ≈ 0.58M pl . At this field excursion, the axion has already traversed M π natural periods. Because back-reaction effects are weak when ψ 1 the inflaton can go through many of its periods with full computational control if M is suitably large 11 .
Comparison with the Wilson-Contour of B 2
We have argued that the integral of C 2 over the 2-cycle of T 1,1 gives rise to an ultra-light mode. Moreover we claim that all other modes are stabilized at least at the warped KKscale m wKK . For this we assume of course that the axio-dilaton is stabilized in the UV by ISD-fluxes on cycles not relevant to the local KS-throat geometry. Perhaps the most obvious candidate that naively seems to be similarly light is the analogous ansatz for B 2 . This however does not give rise to a similarly light mode as can be seen from (48). Integrating this over the throat from the tip up to a position τ 1 (let us call this region C τ 1 ) yields: where N (τ 1 ) is the D3-brane charge integrated over C τ 1 and we have used that ∂C τ 1 = T 1,1 | τ =τ 1 in the second step. Inserting our ansatz (far from the tip) δC 2 ∼ f (τ )ω 2 (with harmonic two form ω 2 ) in (65), the LHS is left unchanged because F 3,KS −→ F 3,KS + dδC 2 and dδC 2 ∼ f (τ )dτ ∧ ω 2 wedges to zero with H 3 of the KS/KT-solution. It therefore does not enter as a source for additional 5-form flux in (48). This can also be seen from (58) which is left unperturbed far from the tip where f ≈ g. 9 In our convention the periodicity of the canonically normalized axion is φc −→ φc + 2πf . 10 By this we mean that the Supergravity approximation is marginally valid at the tip of the throat, gsM ∼ O(1), and the compact CY is of the same size as the UV-region of the throat. 11 Note that the Supergravity approximation becomes better at large M because the typical length scale at the tip is R 2 ∼ gsM α such that the Supergravity approximation is good when both gsM 1 and gs 1.
Choosing the same ansatz for δB 2 instead, there is in contrast a non-trivial wedge-product between dδB 2 and F 3 of the KS-solution such that This can again be directly seen in (48): The five-form flux thus changes linearly in the perturbation also far away from the tip of the throat. It then enters as a source on the RHS of (47), which takes significant (i.e. not exponentially suppressed) values in the whole bulk. This can be interpreted as an effective 5d mass term for the B 2 axion and therefore leads to a 4D-mass of order of the warped KK-scale m wKK .
Conclusion
In this paper we presented a new idea for axion monodromy inflation in which the inflaton is the lightest Kaluza-Klein mode of the RR-2-form potential C 2 wrapped on a homologically trivial 2-cycle. One of the crucial technical points is that the mass of the lightest Kaluza-Klein mode is exponentially lower than that of the next excited mode, thus making this mode an ideal inflaton candidate. The monodromy arises due to the homological triviality of the 2-cycle similar to models proposed in [31], rather than due to a coupling to branes. Consequently, our construction does not require the presence of brane-antibrane pairs, thus avoiding the associated back-reaction issues [34,35]. Crucially, the exponential mass-suppression is due to the S 2 shrinking to zero size only in IR regions. This is why we base our model on the 'double throat' shown in Figure 1. Because back-reaction on other Supergravity fields cannot be neglected at large field excursion the non-linear Supergravity equations that govern the back-reaction are derived. Their numerical evaluation is left for future research.
The full type IIB Supergravity equations also show that the shift-symmetry of C 2 is preserved in the warped background except for the small monodromic potential that is generated at the tip. This in contrast is not true for the analogous Wilson-contour of B 2 .
The perhaps most obvious open question that needs to be addressed by future research is the numerical solution of the equations of motion that were derived. From this the backreacted potential and the maximal controlled field excursion of the model can be inferred.
Naturally then, the question arises what the KS-throat decays to once the field excursion is set beyond its critical value.
Cornering the question of back-reaction from various dual descriptions seems to be a promising path towards gaining analytical insight: In the dual gauge theory the Wilson contour of C 2 at fixed radial position corresponds to the difference of θ-angles of the two gauge group factors. Since θ-angles are left unchanged under the cascade of Seiberg-dualities [42], the nearly constant profile of C 2 seems indeed to correspond to the correct Supergravity dual. It would be interesting to gain analytical insight into the back-reaction of our mode through this dual picture. Furthermore, in a T-dual type IIA picture the Wilson contour of B 2 corresponds to the distance between two NS5 branes with D4-branes suspended between them [18]. An analogous interpretation for the Wilson contour of C 2 that makes the monodromy manifest would be desirable and could also help address the question of back-reaction analytically.
It remains to be seen if the back-reacted potential extracted from the numerical analysis of the type IIB equations of motion is compatible with inflation in general and current observational constraints in particular. A simple quadratic potential is strongly disfavored by the latest data [16]. Back-reaction effects however generically lead to a flattening of the potential [30,64] such that the model may well be in accord with current data.
Overall, we observe that our proposal realizes axion monodromy for a fairly minimal amount of ingredients. Given this relative simplicity and the high level of sophistication with which throat geometries can be controlled [65,66], we expect our model to be a promising arena for further investigations into the viability of large field inflation in string theory. Regardless of the phenomenological implications, we would even like to hope that the possibility of large field inflation could be firmly established based on our simple scenario.
Acknowledgments
We
The 5-form field strength contributes as 1, 1, 1, 1, 1, 1) and the axio-dilaton as g 2 s e 2φ (F while the Ricci-tensor can be computed to give where and There are now seven Einstein equations for the five functions A, B, C, D, E. If we had allowed for a term d(τ )(g 1 g 4 − g 2 g 3 ) in the 10d metric and Θ τ ∝ L(τ )dτ with L(τ ) = D(τ ) there would be seven equations for seven functions. Because such a more general ansatz can always be brought to the form that we have specified by a suitable coordinate re-parametrization the seven Einstein equations are not all independent and we do not expect the resulting system of differential equations to be overdetermined. | 7,170.8 | 2015-12-14T00:00:00.000 | [
"Physics"
] |
Production of ( anti-) ( hyper-) nuclei at LHC energies with ALICE
The ALICE experiment at the LHC has measured a variety of (anti-)(hyper-)nuclei produced in Pb–Pb collisions at √ sNN = 5.02 TeV and at 2.76 TeV. In addition, a large sample of high quality data was collected in pp collisions at √ s = 7 TeV and 13 TeV and in p-Pb collisions at √ sNN = 5.02 TeV. These data are used to study the production of different (anti-)(hyper-)nuclei in the collisions, namely (anti-)deuteron, (anti-)3He, (anti-)alpha and (anti-)3 Λ H. The identification of these (anti-)(hyper-)nuclei is based on the energy loss measurement in the Time Projection Chamber and the velocity measurement in the Time-Of-Flight detector. In addition, the Inner Tracking System is used to distinguish secondary vertices originating from weak decays from the primary vertex. New results on deuteron production as a function of multiplicity in pp, p–Pb and Pb–Pb collisions will be presented, as well as the measurement of 3He in p–Pb and Pb– Pb collisions. Special emphasis will be given to the new results of the (anti-)3 Λ H in its charged-two-body decay mode. The large variety of measurements at different energies and system sizes constrains the production models of light flavour baryon clusters, in particular those based on coalescence and the statistical hadronisation approaches.
Introduction
The nuclei, anti-nuclei and hyper-nuclei formation mechanism in high energy collisions is still unknown.Thanks to the LHC which provides pp, p-Pb and Pb-Pb collisions at the highest energy ever reached in the laboratory, the ALICE experiment is able to shed light on the phenomenology of (anti-)(hyper-)nuclei production.The latest results on the production spectra of (anti-)deuteron and (anti-) 3 He in pp collisions at √ s = 13 TeV and Pb-Pb collisions at √ s NN = 5.02 TeV will be discussed here.The results are compared with the expectation from the statistical hadronisation and the hadron coalescence models.Finally the new measurement of the (anti-) 3 Λ H lifetime by the ALICE experiment, which is the most precise ever performed, will be discussed.
Analyses details
The key features that allow ALICE to measure (anti-)(hyper-)nuclei are the precise vertexing, tracking capabilities and the redundancy of particle identification detectors.
Tracking and vertexing are performed using the Inner Tracking Systems (ITS), a silicon tracker featuring 6 cylindrical layers, and the Time Projection Chamber (TPC) [1].Thanks to the extended e-mail<EMAIL_ADDRESS>and the Lévy-Tsallis in pp [9]) to extrapolate the production yield in the unmeasured p T regions.lever arm and a maximum solenoidal magnetic field of 0.5 T, the momentum resolution is better than 1% for particles with p ≤100 GeV/c [2].Tracks reconstructed with points in the innermost ITS layer have a pointing resolution better than 300 µm [2].For (anti-)(hyper-)nuclei identification the TPC specific energy loss signal (dE/dx) is used to clearly separate nuclei with Z=2 over the full measured momentum range from the bulk of the produced charged particles.Lighter nuclei with Z=1 can be identified clearly by means of the TPC dE/dx only at low momenta (e.g.p ≤ 1.4 GeV/c for (anti-)deuterons), but thanks to the Time Of Flight detector (TOF) it is possible to identify, for instance, (anti-)deuterons up to p T = 6 GeV/c using a statistical unfolding technique [3].In order to identify (anti-) 3 Λ H decaying in (anti-) 3 He and π −(+) it is necessary to reconstruct its decay vertex.For each possible pair of (anti-) 3 He and π −(+) the point of closest approach between the two tracks is computed and a set of topological selection is applied.A pair of (anti-) 3 He and π −(+) that passes all these selections is considered to be a (anti-) 3 Λ H candidate.Finally a two components fit is performed to extract the yield of (anti-) 3 Λ H, similarly to what is illustrated in [4].
Results
The production spectra of deuterons in Pb-Pb collisions at √ s NN = 5.02 TeV and in pp collisions at √ s = 13 TeV are shown in Figure 1 (left).The typical hardening of the spectra with increasing centrality, already observed for lighter particles [5,6] and for deuterons [3,7] at lower energies, is visible here and in the 3 He spectra (right panel of Figure 1): this pattern suggests an increasing radial flow with increasing centrality.Moreover, the ratio between the production spectra of nuclei and their anti-matter partners is compatible with unity [10], like for the other baryon over anti-baryon ratios at the LHC energies.This observation is compatible with expectations from the statistical hadronisation [11] and the hadron coalescence models [12].In both approaches the abundances of a (anti-)nuclide X with mass number A are coupled to those of (anti-)protons according to the following relation: X/X ≈ (p/p) A .This expectation for the coalescence model holds only if the spectra of (anti-)neutrons are compatible with those of (anti-)protons.Under this assumption it is also possible to compute the coalescence parameter B A which is defined for a nuclide with mass number A as the ratio between the nuclide invariant production spectrum and the proton spectrum to the power of A. Figure 2 shows the coalescence parameters for deuterons and 3 He.In the simple hadron coalescence approach the coalescence parameter is assumed to be momentum independent.However, the presence of correlations in the proton and neutron initial emission, such as jets or common emission radius from the thermal source [13], can introduce a momentum dependence similar to the measured B A .
Another observable qualitatively described by the simple hadron coalescence is the deuteron over proton ratio as a function of multiplicity in small systems, see Figure 3.As more protons and neutrons are created in events with increasing multiplicities, the probability of finding two nucleons close in momentum space and consequently the probability of forming a deuteron increase in small systems.However, this scaling breaks in semi-central Pb-Pb collisions where the deuteron over proton ratio as a function of multiplicity becomes flatter.A flat nucleus over proton ratio as a function of multiplicity, observed also in the case of 3 He shown in Figure 3, is expected if the statistical hadronisation is the underlying production mechanism in nucleus nucleus collisions.At very high multiplicity the measurements of the deuteron over proton ratios in Pb-Pb collisions at √ s NN = 2.76 TeV and √ s NN = 5.02 TeV show a hint of suppression.This kind of behaviour might point to the presence of a hadronic rescattering phase that might break nuclei, similar to what is proposed in [14].
Finally, the ALICE collaboration has been able to perform one of the most precise measurements of the 3 Λ H lifetime using the Pb-Pb data sample at √ s NN = 5.02 TeV collected in 2015.The lifetime was measured analysing the 3 Λ H → 3 He + π − and the 3 ΛH → 3 He + π + decays in different ct intervals and performing an exponential fit to the obtained dN/d(ct) distribution, similarly to what has been done in [4].The preliminary measured value τ = 237 +33 −36 (stat.)± 17 (syst.)ps is compatible with both the previously computed world average (τ = 216 +18 −16 ps, [4]) and with the free Λ lifetime.
Conclusions
The harvest of results from the LHC Run 2 data in the (anti-)nuclei sector just started.The new results at the current LHC top energies confirm the picture depicted by the Run 1 results [3,7]: at the LHC energies the nuclei production mechanisms in small systems and in Pb-Pb seem to be different.On one hand the simple hadron coalescence describes qualitatively the measurements of deuteron over proton ratio and B 2 up to peripheral Pb-Pb collisions.On the other hand the steady rise of the B 2 with momentum and the flattening of the deuteron over proton ratio in mid-central and central Pb-Pb collisions require additional assumptions to reconcile the coalescence expectations with ALICE results.The larger data sample expected in the upcoming 2018 Pb-Pb run will help to refine the results for 3 He and possibly extend the measurements to heavier nuclei to test the model predictions with increased precision.However, hadron coalescence might not be the production mechanism of (anti-)nuclei in Pb-Pb collisions: the flat d/p ratio starting from mid-central collisions is in agreement with the expectations of the statistical hadronisation model.Further studies on the systematic uncertainties will help to establish whether the decrease of the d/p ratio in very central Pb-Pb collision is significant.
In the strange nuclear sector, the new ALICE measurement of the 3 Λ H lifetime is compatible with both the current world average and with the theoretical expectation (i.e. the free Λ lifetime [15]).
Figure 1 .
Figure 1.Transverse momentum spectra of deuterons (on the left) and 3 He (on the right) in Pb-Pb collisions at √ s NN = 5.02 TeV.The p T production spectrum of deuterons in pp at √ s = 13 TeV is also reported in the left panel.Statistical uncertainties are represented as vertical bars whereas boxes represent the systematic ones.The dashed lines represent the fit to the individual spectra (Blast Wave function in Pb-Pb[8] and the Lévy-Tsallis in pp[9]) to extrapolate the production yield in the unmeasured p T regions.
Figure 2 .
Figure 2. Coalescence parameters of deuterons (B 2 , left panel) and 3 He (B 3 , right panel) measured in Pb-Pb collisions at √ s NN = 5.02 TeV.The B 2 measured in pp collisions at √ s = 13 TeV is also shown in the left panel.Statistical uncertainties are represented as vertical bars whereas boxes represent the systematic ones.The ALICE preliminary proton p T -production spectra [6] are used to compute the coalescence parameters.
Figure 3 .
Figure 3. Deuteron (left panel) and 3 He (right panel) over proton ratios as a function of the charged particle multiplicity in the laboratory in different collision systems.Statistical uncertainties are represented as vertical lines whereas boxes represent the systematic ones. | 2,298.8 | 2018-02-01T00:00:00.000 | [
"Physics"
] |
Vehicle Counting and Motion Direction Detection Using Microphone Array
—This paper describes a method for counting the vehicles and estimating the direction of their motion using three equidistantly spaced microphones. Paper outlines hardware and software design. Algorithm is based on sound delay estimation between microphones. To implement vehicle counting, sound wave arrival angle is calculated from the derived delays. All measurements and tests took place in real traffic flow under different weather conditions.
I. INTRODUCTION
Increasing number of vehicles demands to develop more sophisticated traffic management systems.To manage the traffic, these systems require different types of information about vehicle (count, speed, line occupation etc.).To gather such information invasive or non-invasive sensors can be used for vehicle detection.Non-invasive sensors application for traffic monitoring and counting are becoming more and more popular in today's traffic organisation and intelligent traffic systems.Main advantages over the invasive sensors are: easier installation and maintenance without disturbing the flow of traffic.
One of the ways of gathering information about traffic flow is to record sound of moving vehicles by applying microphone arrays and use signal processing to extract information about vehicle.Microphone arrays belong to non-invasive sensor group.Acoustic sensor can be a single microphone for detecting the emergency vehicle or the arrays of multiple microphones can be used for vehicle motion direction detection, counting and classification.Roads are noisy environment with different types of noise sources (car tires, engines, different types of car moving parts and wind noises).One of the advantages of microphone arrays, compared to the video registration, is that they can deliver reasonable performance in night and in limited visibility conditions.
Moving vehicle noise depends on vehicle speed.Most of noises coming from moving vehicles are noises from tires, noises generated by vehicle aerodynamics and noises Other noise strongly depends on the vehicle type.Vehicles of the same type under identical condition are emitting similar acoustic signals.By changing vehicle speed and road conditions the emitted signals will differ.There is also difference in the emitted sound from the heavy and light vehicles.The microphone array receives a mixture of vehicle acoustic emissions and ambient noise.Ambient noise (for example wind) can be filtered by hi-pass filter.It is possible to extract spectral features from the signals recorded from moving vehicles.These spectral features describe vehicle properties like: motion direction, count, also classification (heavy, lightweight, emergency vehicle etc.).
II. RELATED WORK
There are many articles and research in the field of microphone array signal processing with applications in speech recognition, speech extraction applying acoustic beam-forming and cross power spectrum phases analysis algorithm [1], [2].
The use of microphone array or single microphone in traffic management is not a new idea.There are many studies of microphone arrays applications for vehicle detection [3]- [6].Many articles examine single microphone or microphone array application of the vehicle detection and counting.For example paper [7] describes tests in Tintersection with the microphone array consisting of six microphones.The detection results obtained with weighted sum of cross-correlation functions are compared by two conventional methods used for microphone array signal processing.The article [8] describes different vehicle emitted sound parameterization methods which can be used in classification systems.Some articles describe the road condition detection, like frozen, dry or wet using single microphone.Article [9] describes experiments with single microphone for ice detection on the road surface.An interesting method is the microphone placement on the moving vehicle close to the wheel.Also emergency transport presence and approaching direction detection from sirens can be obtained using single microphone or microphone array if necessary [10], [11].All of the mentioned applications have similar approaches based on the sound delay estimation.
III. VEHICLE COUNTING ALGORITHM
Proposed vehicle counting algorithm consists of several steps which are illustrated in Fig. 1.The first step is to find delays between microphones, which is the most important part of the vehicle counting algorithm.Delays are estimated in relation to the reference channel by applying (Generalized cross-correlation (GCC)) method [12].The second microphone is chosen as a reference channel.In calculations far field assumption is applied.We assume that propagating sound waves are plane waves, neglecting curvaceous sound wave properties.This assumption simplifies the estimation of a car sound arrival angle by applying relatively simple geometric calculations (Fig. 2) using (1) based on estimated delays, where θ is angle in radians, τ is estimated delay in seconds, c -acoustic velocity (343 m/s) and d -distance between microphones After calculating the absolute values of sound arrival angle, the peaks which represent the vehicle presence are found.To filter calculated angle values threshold is applied.Threshold value is 70 degree, it was found experimentally.In further calculations only components above 70 degree are used (all components under 70 degree are removed).To prevent false peaks appearance above 70 degree we use nearest peaks neighbour analysis.One angle value calculation is performed on 250 ms of signal length (window size).The window moving step is equal to window length.In this time interval vehicle travels 6 meter distance.If motion speed is 90 km/h, that corresponds from 0 -20 degree angle deviation.This angle deviation depends on the location of vehicle in the relation to the microphone array.If two closest neighbour angle values differ more than deviation of 20 degree we conclude that it is not the vehicle.After this operation algorithm calculates motion direction.
The first step of vehicle motion direction estimation is to calculate first order derivative from the obtained delays.
First order derivative represents the rate of change of the delays.From the values of the low and high levels of the first order derivative vehicle motion direction is estimated.Low and high thresholds are empirically obtained.
IV. HARDWARE DESIGN AND IMPLEMENTATION
To record the sound emitted by the vehicle data collecting device were designed.Device consists of three WM7110 MEMS (Micro-electromechanical systems) analogue microphones, amplifier and hi-pass filter.The first tests were performed using analogue capacitive microphones but because of the large parameter dispersion these sensors were not used in the final prototype.Analogue filter is used to filter frequencies under 400 Hz.This frequency band is chosen because vehicles radiate nose above this frequency but most of the wind noise components are below this frequency.In our prototype microphones are equidistantly spaced with 10 cm interval (Fig. 2: d1 and d2).This distance depends on sound wavelength.AD1974 ADC is used to process analogue signals with 48 kHz sampling frequency.ADC is connected to ALTERA Cyclone III board via I2S interface.All data is sent to PC for further processing via USB interface.
Sound signal processing algorithms are developed using Matlab.Designed system block diagram is illustrated in Fig. 3.In this development stage hardware architecture provides only non -real time implementation.The second prototype is planned as standalone device including vehicle counting, direction detection and vehicle classification features.
V. MEASUREMENTS AND EXPERIMENTAL RESULTS
All sound recordings of the moving vehicles were performed on two lane highway approximately 400 meters from traffic light.Microphone array was located parallel to the road about 6 m from first lane and 1 m above the ground.Flow of the vehicles is grouped because of the regulated intersection.Designed microphone array prototype is displayed in Fig. 4. Approximately two hours of traffic flow were recorded under different weather conditions (rain, windy, wet pavement, sunny).A recording mostly consists of sounds from heavy vehicles with trailers (about 65%).The average speed of all vehicles is about 70-100 km/h.In addition to the sound, video recordings were made to capture vehicle movement for comparison purposes.Video camera was fixed above microphone array on a tripod.
Figure 5 illustrates the sound signal that was emitted by the vehicles and recorded with one of the three microphones.This test was performed in sunny day in low wind and dry pavement condition.Average speeds of the vehicles were about 90 km/h (25 meter per second).Figure 5 displays a signal filtered using 0,4 kHz hi-pass analog filter.In this case displayed sound signal is emitted from nine vehicles.From the graph with the recorded sounds, it can be seen that the vehicles moving near the microphone from 25 to 35 second can't be detected without data processing.There are four moving vehicles at this time interval.The following pictures from 6 to 10 are signal processing results obtained from signal displayed in Fig. 5.
Figure 6 illustrates the spectral estimate of the sound emitted by the vehicles.The spectral estimate is obtained using Short time Fourier transform (STFT).Figure 6 displays that the most of the spectral components are in frequency range from 0,4 to 3 kHz.As mentioned earlier, before analog signal sampling first order hi-pass filter is used.A small part of signal spectral energy emitted by the vehicles is above 3 kHz.Using ideal FFT filter all spectral components above 5 kHz were filtered and using IFFT ideally filtered signal spectrum were transformed back to time domain.That was done for experimental purposes to test which of the spectral components are important for vehicle detection.When delays and angle of the sound arrival was calculated random peaks appeared in result which generates false positive error in vehicle detection.Identical experiment using ideal FFT low pass filter were performed at 10 kHz cutoff frequency.
This experiment using ideal FFT filter with different cutoff frequencies showed that vehicle emitted sound spectral components over 3 kHz which energy is small compared to components below 3 kHz are significant for sound arrival angle detection.This aspect is important for choosing sampling frequency.In our experiments we tested different sampling frequencies: 10 kHz, 15 kHz, 24 kHz and 48 kHz.The best results are achieved using 48 kHz sampling frequency.
Figure 7 displays signal energy emitted by the vehicle that can also be used to detect vehicle presence.Disadvantage of this method is poor vehicle detection when object distance is small.An experiment shows that this distance is in range of 10 meters.The video recording shows that in time interval from 25 to 35 second four vehicles passed by microphone array but only three of them were detected using spectral energy estimation (Fig. 7).
Proposed method for vehicle counting by estimating angle of arrival in time interval from 25 to 35 second finds all four vehicles (Fig. 10).The analyzed one minute signal (Fig. 5) contains (ground truth) nine vehicles.In this case our method discovers all nine vehicles.If random noise is present it is filtered by applying closest neighbour angle elimination method, described in third section of paper.
The second part of algorithm describes motion direction detection.Figure 8 displays acoustic signal (Fig. 5) delays from two microphones.The delays are obtained using GCC method [12].The first order derivative is applied to calculate slope of delays.The direction of motion is obtained by comparing the peaks of calculated delay derivatives.The first vehicles direction is obtained by the value of delay.Positive delay means the car is approaching from the right direction, negative -from left.If the first peak of derivative is smaller than the second then vehicle is approaching from the same direction as the first vehicle.Figure 10 illustrates vehicles sound arrival angle in relation to microphone array.Circled points show detected vehicles.
VI. CONCLUSIONS
This paper proposes the real traffic detection method using three equidistantly spaced microphones.Proposed method includes vehicles counting and motion direction detection.The method is based on the sound delay estimation of the moving vehicle.Sound delay and sound waves angle of arrival are estimated by using three microphones.
Tests were carried out on the highway outside the city and other populated areas.The tests were made in different weather conditions, except snowing.About 2 hours of recordings were analysed.Vehicle counting results are summarized in Table I.In total 871 vehicles (ground truth) were counted based on the analysed video recordings of 2 hours.False positive detection results occur in strong wind weather conditions; also wet pavement is the source of false positive result.False negative result detection was observed when multiple vehicles pass in front of the microphone array.Situations such overtaking, small distance withholding are also false negative detection reasons.In case when vehicle speed is approximately 90 km/h the smallest distance between vehicles for counting algorithm is 7 -10 m.
The accuracy of motion direction is 65 percent.The precision was calculated based on recorded video.Motion detection accuracy is lower than vehicle counting accuracy because distances between vehicles in this case have to be longer than 10 meters (in case of light vehicles).Because all traffic data where acquired approximately 400 m from regulated intersection, distance between vehicles in many cases are smaller than 10 meters.
This approach showed that heavy vehicle emitted noise suppresses immediately followed light vehicle sound if the distance between objects is smaller than 10 meters.This distance between heavy and light vehicles strongly depends on window size taken to compute angle of sound arrival.If two cars approaches from both sides and meets in front of microphone array only one car is detected.
Fig. 1 .
Fig. 1.Block diagram of vehicle counting and direction detection algorithm.
Fig. 3 .
Fig. 3. Block diagram of designed system for vehicle detection.
Fig. 5 .
Fig. 5.One minute recording of highway traffic sound with two lanes.
Figure 9
Figure 9 illustrates estimated delays and first order derivative of one minute traffic flow on highway.
TABLE I .
VEHICLE DETECTION RESULTS. | 3,004.4 | 2013-07-10T00:00:00.000 | [
"Computer Science",
"Engineering"
] |
Investing in mental health and well-being: findings from the DataPrev project
A systematic review was conducted to determine the extent to which an economic case has been made in high-income countries for investment in interventions to promote mental health and well-being. We focused on areas of interest to the DataPrev project: early years and parenting interventions, actions set in schools and workplaces and measures targeted at older people. Economic evaluations had to have some focus on promotion of mental health and well-being and/or primary prevention of poor mental health through health-related means. Studies preventing exacerbations in existing mental health problems were excluded, with the exception of support for parents with mental health problems, which might indirectly affect the mental health of their children. Overall 47 studies were identified. There was considerable variability in their quality, with a variety of outcome measures and different perspectives: societal, public purse, employer or health system used, making policy comparisons difficult. Caution must therefore be exercised in interpreting results, but the case for investment in parenting and health visitor-related programmes appears most strong, especially when impacts beyond the health sector are taken into account. In the workplace an economic return on investment in a number of comprehensive workplace health promotion programmes and stress management projects (largely in the USA) was reported, while group-based exercise and psychosocial interventions are of potential benefit to older people. Many gaps remain; a key first step would be to make more use of the existence evidence base on effectiveness and model mid- to long-term costs and benefits of action in different contexts and settings.
withdrawal from the workforce on health grounds (McDaid, 2011).
While these are serious impacts, they are in themselves insufficient to justify investment in measures to promote mental health and wellbeing. For this, it is important not only to identify robust evidence-informed actions, but also to look at their costs and resource consequences, within and beyond the health system. Resources are always finite, with many potential alternative uses, and careful choices have to be made on investment and priority setting. It is perhaps even more critical to highlight whether investment in the promotion of mental health and well-being might represent good value for money and help avoid future costs of poor mental health during the current austere climate when health and other public sector budgets are under substantial pressure, and when mental health promotion may not be seen as a high priority for policy makers .
As part of the EC funded DataPrev project, a systematic review was conducted to identify the state of the evidence base on the use of economic evidence in helping to make the case for investment in mental health and well-being in the four areas of focus to the project: early years and parenting interventions, actions set in schools and workplaces and measures targeted at older people.
METHODS
Our objective was to identify economic evaluations, i.e. studies comparing the effectiveness and costs of two or more health-focused interventions, to promote mental health and wellbeing and/or prevent the onset of mental health problems.
Inclusion and exclusion criteria
Two distinct types of study were eligible for inclusion. First, economic evaluations conducted concurrently or retrospectively alongside a randomized controlled trial. An exception to this criterion was applied to workplace health promotion interventions where controlled trials are rare; in this case other empirical study designs alongside an economic analysis were also eligible. Economic evaluations conducted using a modelling approach, whereby effectiveness data were collected from one or more previous controlled studies and then combined with data on costs, were also included. Economic evaluations had to be consistent with different approaches commonly applied in health economics, including cost-effectiveness, cost -benefit, costconsequence, cost -utility and cost-offset analyses. While we cannot discuss the differences between these approaches here, the interested reader can refer to numerous guides, e.g. (Drummond et al., 2005;Shemilt et al., 2010).
To be eligible for inclusion studies also needed to include either a measure of positive mental health, e.g. use of the SF-36 mental health summary scale or other measures of quality of life, specific measures of well-being or alternatively quantify the prevention of psychosocial stress and/or mental disorders. We excluded studies relating to the prevention of dementia, as well as those focused on individuals with learning difficulties from our analyses. Interventions needed to have a primary objective of promoting health. This meant that we excluded some education and child care centred interventions that had subsequently been shown to have a positive impact on mental health (among other outcomes) (Barnett, 1998;Barnett and Masse, 2007).
Papers that focused on the treatment of individuals with existing mental health problems were excluded, with the exception of studies that looked at how the treatment of parents with mental health problems might promote/ protect the mental health of their children, as well as those reporting proxy outcomes, such as improvements in parent-child interaction and the prevention of child abuse. Children were assumed to be between the ages of 0 and 16, while studies in respect of older people focused on people aged 65 plus.
Search process
A search strategy designed to identify economic evaluations in bibliographic databases (Sassi et al., 2002) was combined with a range of mental health promotion/mental disorder terms and a set of population/setting specific keywords and phrases. Mental health-related terms and concepts included in the search included mental health, positive mental health, mental and emotional well-being, personal satisfaction, quality of life, happiness, resilience, energy and vitality. Health promotion and prevention-Investing in mental health and well-being i109 related keywords and phrases were also combined with terms related to poor mental health, including psychological stress, post-natal/postpartum depression, conduct disorder and child behavioural disorders.
We searched PubMed, PsycINFO, EMBASE, CINAHL, PAIS, Criminal Justice Abstracts, Web of Science, Scopus, EconLit and the National Health Service (NHS) Economic Evaluation Database at the University of York. Only results that reported abstracts (or chapter summaries) in English were included; geographical coverage was limited to the European Economic Area, plus EU Candidate Countries, Switzerland and other Organisation for Economic Co-operation and Development (OECD) members. Our review covered the period from January 1990 to December 2010. The electronic search was complemented by a limited search for key terms in Google Scholar, the general Google search engine and scrutiny of relevant websites, e.g. think tanks, universities, government departments and agencies. We also undertook a handsearch of a small number of journals and examined the reference lists of included studies, as well as citations of papers that met our inclusion criteria.
In addition, we also looked for any economic analyses of mental health promoting interventions previously shown in companion systematic reviews on effectiveness conducted as part of the DataPrev study to be effective in promoting mental health and well-being. Where these reviews identified evidence of the impact of an intervention on mental health and well-being, any studies that looked at the economic case for investment in those interventions, even if focused on non-health benefits, such as improved educational attainment, reduced crime and violence, were then eligible for inclusion.
References were initially screened independently by two reviewers (D.M. and A.P.) on the basis of their abstracts/summaries to determine whether they met study inclusion criteria. In the case of disagreement the two reviewers discussed the paper and came to a final decision on inclusion/exclusion, erring on the side of inclusion where no easy agreement could be reached. The full text of all references appearing to meet initial inclusion criteria was then retrieved and a final assessment made. Ultimately included studies were coded and stored in an Endnote database. An assessment of the quality of studies was also made, making use of two published economic evaluation checklists (Drummond and Jefferson, 1996;Evers et al., 2005). Overall this process meant that .3000 references were assessed (see Figure 1).
Parenting, early years and school-based interventions
There has been a considerable body of research into the effectiveness of interventions to promote/protect the mental health and wellbeing of children and their parents, both within and external to school settings (Adi et al., 2007a, b;Dretzke et al., 2009); there is also a small but growing number of studies looking at the economic case for taking action, albeit largely set in either a USA or UK context. We also identified one study protocol for an economic evaluation of an internet-based group intervention to prevent mental health problems in Dutch children whose parents have mental health or substance abuse problems (Woolderink et al., 2010). Overall the results are mixed, as the summary of findings from 26 papers and reports in Tables 1 and 2 indicate. Table 1 includes several studies looking at the impact of health visitors, including the wellcited Nurse Family Partnership programme developed in New York in the 1980s (Olds et al., 1993). Focusing on new mothers, but with a special emphasis on teenage, single-and lowincome mothers, the study followed 400 mothers and their children over a 15-year period. Looking at a broad range of outcomes going beyond positive maternal and child mental health outcomes, an initial analysis reported net costs per woman of $1582 (1980 prices) over the first 4 years for the whole population, but net savings of $180 per high-risk woman (Olds et al., 1993).
Empirical studies
Home visiting programmes have also been examined in England; some focused directly on child mental well-being, others on avoiding post-natal depression, a risk factor for poor child mental health (Murray, 2009). A controlled trial of an intensive home visiting programme and social support programme for i110 D. McDaid and A-La Park vulnerable families where children could be at risk of abuse or neglect reported a cost per unit improvement in maternal sensitivity and infant cooperativeness of £3246 (2004 prices) (Barlow et al., 2007;McIntosh et al., 2009). The challenge with such a finding, however, is judging whether this well-being improvement represents value for money, as it uses a clinical outcome measure which cannot be compared with other uses of resources within the health-care system. Both cost -utility analyses where outcomes are measured in a common metric, such as the Quality Adjusted Life Year (QALY) where a maximum cost per QALY deemed to be cost effective can be determined in different contexts, or cost -benefit analyses where both outcomes and costs are measured in monetary terms can be used to overcome this problem, although neither approach is without its own limitations (Kilian et al., 2010).
In England, a randomized controlled trial of health visitor delivered psychological therapies for women at high risk of post-natal depression improved outcomes at lower costs than health visitor usual care. There was a 90% chance that the cost per QALY gained would be ,£30 000; a level generally considered to be cost effective in an English context (Morrell et al., 2009). Another trial of women at high risk of post-natal depression compared health visitor delivered counselling and support for motherinfant relationships to routine primary care, finding that if society was willing to spend £1000 to prevent 1 month of post-natal depression then the intervention would have a 71% chance of being cost effective with mean net benefits of £384 (2000 prices) (Petrou et al., 2006). This contrasted with an earlier study on the use of post-natal support workers to reduce the risk of post-natal depression which did not appear cost effective (Morrell et al., 2000). However, the former study needs to be interpreted carefully as neither the change in costs or outcomes in the trial were significant and a comparable measure such as the QALY was not be used. Covering a longer time period and looking at additional benefits to children and mothers may have strengthened study findings.
Compared with standard health visitor care, no effectiveness or economic benefit was found in making use of supportive home visits to ethnically diverse mothers in London (Wiggins et al., 2004(Wiggins et al., , 2005. Home visiting was also compared with participation in a mother -child attachment group intervention in Canada. While no difference in effects was reported, costs were significantly lower in the attachment group (Niccols, 2008). We also found a recent Investing in mental health and well-being i111 (Foster, 2010)
Continued
Investing in mental health and well-being i115
Continued
Investing in mental health and well-being i119 The cost of implementing Triple P to one cohort of 2 year olds would be AUD 9.6 million. The average cost per child in the cohort would be AUD 51 RCT, randomized controlled trial; CBA, cost -benefit analysis; CEA, cost-effectiveness analysis; CCA, cost-consequences analysis; CUA, cost -utility analysis; COA, cost-offset analysis.
Australian study that reported that the provision of advice and materials within a maternal and child health centre to mothers of infants with sleep problems had similar costs but better mental health outcomes for mothers and improved sleep patterns for infants compared with standard clinic consultations (Hiscock et al., 2007). As Table 1 indicates, a number of economic evaluations of parenting studies conducted alongside randomized controlled trials have been published, some set in schools, others focused on pre-school age children. In addition we identified one published study protocol for an ongoing evaluation in Wales (Simkiss et al., 2010). An evaluation of the Webster-Stratton Incredible Years parenting programme in Wales, while finding the intervention to be costeffective for all 3-5-year-old children at risk of conduct disorder, suggested that the intervention would be most cost-effective for children with the highest risk of developing conduct disorder (Edwards et al., 2007). Analysis from a trial looking at 3 -8-year-old children in the USA also suggests that combining the parenting component of Incredible Years with child-based training and teacher training, even though more expensive, can be more cost-effective .
As with many health promotion interventions, benefits are only achieved if there is uptake and continued engagement with an intervention over a period of time. One Canadian study looked at community group versus clinic-based individual parenting programmes; while both approaches were effective in reducing the risk of conduct disorders the community group approach was six times more cost-effective because it reached a larger number of parents (Cunningham et al., 1995). A trial of the Incredible Years Programme, combined with a manualized intervention using reading to promote interaction between disadvantaged parents and their children in London, would however only be cost-effective if uptake and engagement rates could be improved (Scott et al., 2010).
The most negative studies were linked to empirical analysis of the Fast Track programme, a 10-year, multi-component prevention programme implemented in four areas in the USA and focused in part on the promotion of better mental well-being and the prevention of antisocial behaviour and violence. Although this included as one component a school curriculum approach based on PATHS (Promoting Alternative Thinking Strategies), it did not appear to be cost-effective. This may have been partly due to limitations in outcomes data in the study, but even if the intervention could be targeted solely at high-risk children it would only be cost-effective if society was willing to pay more than $750 000 (2004 prices) per case of conduct disorder averted (Foster andJones, 2006, 2007;Foster, 2010). In all of these Fast Track studies no specific monetary valuation was placed on the maintenance of better mental health and well-being, but rather on the longterm consequences to non-health sectors, such as criminal justice.
Modelling studies
As Table 2 indicates, economic models have been used to estimate some of the long-term potential costs and benefits associated with parenting, early years and school-based interventions. Further economic analysis, drawing on 15-year outcome data (Olds et al., 1997) suggested that the economic case for home visiting for all women was much stronger, given the impacts it had in terms of reducing abuse, violence, the need for social welfare benefits and improved employment prospects (Karoly et al., 1998(Karoly et al., , 2005. Benefits outweighed costs by a factor of 5.7 to 1 for high-risk women and 1.26 to 1 for low-risk women. As part of a wide-ranging economic analysis of early intervention programmes commissioned by the Washington State Legislature, several programmes relevant to DataPrev were modelled. It should be noted that the authors of these analyses acknowledged that a limitation of their modelling analysis was that it did not put a monetary value on the economic benefits associated with gains in social and emotional mental well-being or broad health benefits. This was due to the terms of the reference received from the Washington State Legislature, which limited the outcomes for all evaluations to crime, substance abuse, educational outcomes, teenage pregnancy, teenage suicide attempts, child abuse, neglect and domestic violence (Aos et al., 2004).
Nonetheless this Washington State review included further evidence of an economic case for action. Analysis of the Nurse Family Partnership, making use of further updated cost i122 D. McDaid and A-La Park data (Olds et al., 2002) reported a benefit to cost ratio of 2.88 to 1 when modelling benefits to child school leaving age, with major benefits due to crime avoided (Aos et al., 2004). Combining data from several similar home visiting programmes a benefit: cost ratio for programmes targeting high-risk mothers had a 2:1 return on investment, with net benefits per mother of $6077 (2003 prices). (Aos et al., 2004).
Turning to school-based interventions, the Caring School Community scheme developed in the USA (Battistich et al., 1996) and now being implemented in Europe, can be delivered at a cost of $16 per pupil over 2 years, and potentially generate a return on investment of 28:1, even when just looking only at benefits of reduced drug and alcohol problems alone (Aos et al., 2004). Using data from the Seattle Social Development Project, which implemented a teacher and parent intervention including child social and emotional development for 6 years and then followed up these children from age 12 to 21 (Hawkins et al., 2005), costs of $4590 (2003 prices) per child were outweighed by benefits that were three times as great. Again this analysis may be conservative, as no monetary value was placed on the significant improvements seen in mental and emotional health (Aos et al., 2004).
Another school-based intervention that has been modelled is the Good Behaviour Game (GBG), an approach which seeks to instil positive behaviours in children through participation in a game, with prizes given to winning teams who behave better. Potential net cost savings of between $15 and $20 million might be achieved for a hypothetical cohort of 5-and 6-year-old children if the programme could achieve a 5% reduction in special education placements, a 2% reduction in involvement with prison services and a 4% reduction in lifetime prevalence of tobacco use (Embry, 2002). Focusing solely on the economic benefits from evidence on a reduction in tobacco use rather than on any of its mental well-being benefits (Kellam and Anthony, 1998), another analysis of the GBG reported a return of investment of 25:1 (Aos et al., 2004).
As Table 2 shows several economic models have looked at the case for investing in the multi-component, manualized multi-level Triple P-Positive Parenting Programme in a number of different settings. Modelling the potential benefits of universal application of Triple P to the Queensland child population aged 2 -12, the average cost per child would be AUD 34 (2003 prices). It would appear to offer very good value for money when assumed to reduce the prevalence of conduct disorder by up to 4%, generating cost-savings of AUD 6 million. The intervention would have better outcomes and costs would be outweighed by conduct disorder averted as long as the prevalence of conduct disorder was at least 7% Mihalopoulos et al., 2007). In a USA context, an economic model predicted that the costs of Triple P could be recovered in 1 year through a modest 10% reduction in the rate of child abuse and neglect (Foster et al., 2008).
In England, modelling work for NICE (National Institute for Health and Clinical Excellence) looking at the universal use of a teacher delivered PATHS programme for children combined with parent training was reported to have a 66% chance of having a cost per QALY gained of ,£30 000. Combining emotional and cognitive benefits in the model's base case scenario the cost per QALY gained would be £5500 (McCabe, 2008). Other modelling work looking at universal use of social and emotional learning interventions for 11-16-year-old children, and drawing on a review of effectiveness evidence on its application to the prevention of bullying (Evers et al., 2007), suggested that if the intervention reduces victimization by 15% then it would have an 92% of having a cost per QALY ,£30 000 (Hummel et al., 2009).
Promoting mental health at the workplace A number of reviews have looked at evaluations of the effectiveness of interventions delivered in the workplace to promote better mental health and well-being (Kuoppala et al., 2008;Corbiere et al., 2009;Martin et al., 2009a). Actions can be implemented at both an organizational level within the workplace and targeted at specific individuals. The former includes measures to promote awareness of the importance of mental health and well-being at work for managers, risk management for stress and poor mental health, for instance looking at job content, working conditions, terms of employment, social relations at work, modifications to physical working environment, flexible working hours, improved employer -employee Investing in mental health and well-being i123 communication and opportunities for career progression. Actions targeted at individuals can include modifying workloads, providing cognitive behavioural therapy, relaxation and meditation training, time management training, exercise programmes, journaling, biofeedback and goal setting. Tables 3 and 4 summarize key findings on the economic case for investment in workplace mental health promotion from empirical and modelling-based studies. While the costs to business and to the economy in general of dealing poor mental health identified at work have been the focus of attention by policy makers in Europe and elsewhere in recent years (Dewa et al., 2007;McDaid, 2007), less attention has been given to evaluating the economic costs and benefits of promoting positive mental health in the workplace. A recent review for NICE found no economic studies looking specifically at mental well-being at work had been published since 1990 (National Institute for Health and Clinical Excellence, 2009a).
In part this may be due to a lack of incentives for business to undertake such evaluations, as well as issues of commercial sensitivity. There have been few controlled trials of organizational workplace health promoting interventions, let alone interventions where mental health components can be identified and even fewer where information on the costs and consequences of the intervention are provided (Corbiere et al., 2009). Moreover, many actions within the corporate world do not tend to be published in academic journals or books but rather in company literature. This makes studies more difficult to find and a full search of company literature was beyond the scope of our review. Most workplace health promotion evaluations related to mental health have focused on helping individuals already identified as having a mental health problem remain, enter or return to employment (Lo Sasso et al., 2006;Wang et al., 2006;Brouwers et al., 2007;McDaid, 2007;Zechmeister et al., 2008).
In fact, we were able to identify several economic analyses with some focus on mental health promotion (Table 3), largely from a US context where employers have had an not inconsiderable incentive to invest in workplace health promotion programmes, given that they typically have to pay health-care insurance premiums for their employees (Dewa et al., 2007). At an organizational level, modelling work undertaken as part of the UK Foresight study on Mental Capital and Well-being suggests that substantial economic benefits that could arise from investment in stress and well-being audits, better integration of occupational and primary health-care systems and an extension in flexible working hours arrangements (Foresight Mental Capital and Wellbeing Project, 2008). Modelling analysis of a comprehensive approach to promote mental well-being at work, quantifying some of the business case benefits of improved productivity and reduced absenteeism was also produced as part of guidance developed by NICE (Table 4). It suggested that productivity losses to employers as a result of undue stress and poor mental health could fall by 30%; for a 1000 employee company there would be a net reduction in costs in excess of E300 000 (National Institute for Health and Clinical Excellence, 2009b). Another analysis looking at the English NHS workforce reported potential economic gains from reducing absence levels down to levels seen in the private sector that would be equivalent to .15 000 additional staff being available every day to treat patients. This would amount to an annual cost saving to the English NHS of £500 million per annum (Boorman, 2009).
Most analyses have focused on actions targeted at individuals, such as stress management programmes, which are less complex to evaluate. There have been a number of economic assessments of general health promotion and wellness programmes (Pelletier, 1996(Pelletier, , 2001(Pelletier, , 2005(Pelletier, , 2009Chapman, 2005), but few have specifically mentioned mental well-being orientated components, and even when they do include these components they may not report mental health or even stress-specific outcomes. The Johnson and Johnson wellness programme, which includes stress management, has been associated with a reduction in health-care costs of $225 per employee per annum (Ozminkowski et al., 2002), while a 4-year analysis of the Highmark company wellness programme, including stress management classes and online stress management advice, reported a return on every $1 invested of $1.65 when looking at the impact on health-care costs (Naydeck et al., 2008). Neither analysis reported specific impacts on mental well-being or stress. Another study of an intervention to help cope with stress in the computer industry did not find any significant difference in stress levels, but it was associated i124 D. McDaid and A-La Park No ratio reported, as no significant difference in stress, anxiety and coping C: Self-help groups with e-mail personal feedback (partial intervention) and waiting list control 12 months CCA Costs would be lower at $47.50 if delivered by in house medical professionals There was a nearly significant difference in self-reported days of illness for the intervention group But significant 34% reduction in health-care utilization by intervention participants compared with the control groups (p ¼ 0.04) Concluded that this reduction in costs would more than cover the costs of delivering the intervention if delivered by in-house professionals (Renaud et al., 2008), Canada I: Comprehensive health promotion programmes to provide employees with information and support for risk factor reduction, using a personalized approach and involving the organization's management as both programme participants and promoters. Programme includes modules on stress management, healthy eating and physical activity 270 company employees Before and after study. No controls. COA Cost of the intervention not reported Costs avoided not directly reported in monetary terms, but in terms of absenteeism and staff turnover Significant reduction in stress levels away from work as reported using Global Health Profile Score over 3 years falling from 27 to 17% (p , 0.0001). There was also a reduction in feelings of depression with 54.8% of participants stating that they rarely felt depressed after 3 years compared with 38.5% at baseline (p , 0.0001). There was also a reduction in the number of people experiencing signs of stress (p , 0.0001)
Company perspective
No ratio. Significant reduction in high levels of stress, signs of stress and feelings of depression
Continued
Investing in mental health and well-being i127 with a significant reduction in overall reported illness and a one-third decrease in the use of health-care services which would more than cover the costs of the intervention (Rahe et al., 2002). One study that did report mental health outcomes looked at the economic case for investing in multi-component workplace-based health promotion programme ( personalized health and well-being information and advice; health-risk appraisal questionnaire, access to a tailored health improvement web portal, wellness literature, and seminars and workshops focused on identified wellness issues). Using a pre-post test study design, participants were found to have significantly reduced health risks, including work-related stress and depression, reduced absenteeism and improved workplace performance. The cost of the intervention to the company was £70 per employee; there was a 6-fold return on investment due to a reduction in absenteeism and improvements in workplace productivity (Mills et al., 2007).
The experience of employees in another health promotion scheme over 3 years was compared with matched controls. Overall levels of risk to health were significantly reduced, while there was also a significant reduction in the prevalence of depression, although rates of anxiety significantly increased. There were net cost savings from a health-care payer perspective, although the costs of participation in the health promotion programme were not reported (Loeppke et al., 2008). In Canada, an uncontrolled evaluation of a comprehensive workplace health promotion programme, including information for stress management reported a significant reduction in stress levels, signs of stress and feelings of depression at the end of a 3-year study period. While costs of the programme were not reported, staff turnover and absenteeism decreased substantially (Renaud et al., 2008). A small controlled study looking at a programme to prevent stress and poor health in correctional officers working in a youth detention facility in the USA, reported incremental cost savings of more than $1000 over 3 months, although the sample size was too small to be significant. However, the study did not monetize the value of reported productivity gains, while there were positive changes in outlook, attitudes, anger and fatigue (McCraty et al., 2009).
Studies can also be identified where no impacts on absenteeism rates of stress management interventions were identified (van Rhenen et al., 2007). In other cases analyses of a combination of organizational and individual stress management measures did report improvements in emotional well-being, as well as in productivity and reduced absenteeism, but no cost data were provided (Munz et al., 2001). We also identified an ongoing cost -benefit analysis currently being conducted alongside a randomized controlled trial of a mental health promotion intervention to prevent depression targeted at managers in small and medium size companies involving cognitive behavioural therapy and delivered by DVD in Australia (Martin et al., 2009b).
Investing in the mental health and well-being of older people
The final area we reviewed concerned the mental health and well-being of older people. Sixteen per cent of older people may have depression and related disorders; potentially the prevention of such depression, particularly among high-risk groups such as the bereaved, might help avoid significant costs to families, and health and social care systems (Smit et al., 2006). Evaluations from a wider range of countries were identified, most notably from the Netherlands (Table 5). In addition to published studies discussed below, we also were able to identify some ongoing cost-effectiveness studies where protocols had already been published in open access journals (Joling et al., 2008;Pot et al., 2008).
Several studies looked at different types of home visiting interventions to promote wellbeing and reduce the risk of depression, with mixed results. Neither a home visit programme by nurses in the Netherlands nor a programme to promote the befriending of older people in England was found to be effective or costeffective (Bouman et al., 2008a, b;Charlesworth et al., 2008;Wilson et al., 2009). We did identify a cost -utility analysis from the Netherlands conducted alongside a randomized controlled trial comparing a home visiting service provided by trained volunteers with a brochure providing information on depression . It targeted older people who had been widowed for between 6 and 9 months and who were experiencing some degree of loneliness.
Investing in mental health and well-being i129
Continued
Investing in mental health and well-being i131 Although improvements in quality of life were marginal, because of health service costs avoided the intervention had a 70% chance of being cost-effective, with a baseline cost per QALY gained of E6827 (2003 prices). In Canada a home nursing programme used to bolster personal resources and environmental supports of older people was also associated with a reduction in the risk of depression at no additional cost (Markle-Reid et al., 2006). Recently a controlled trial of a stepped care approach for the prevention of depression in older people in the Netherlands was also found to be highly cost-effective at E4367 per depression/anxiety-free year gained (2007 prices) ( Van't Veer-Tazelaar et al., 2010). Economic analyses also supported investment in some different types of group activities. Regular participation in exercise classes by older people was found to have some mental health benefits and be cost-effective from a health system perspective in England with a cost per QALY gained of E17 172 (2004 prices) (Munro et al., 2004). Several studies also reported the beneficial effects to mental health of Tai Chi (MacFarlane et al., 2005), but no formal cost-effectiveness analysis appears to have been undertaken. A study of 166 people randomized to participation in a choral singing group or no action was associated with a reduction in loneliness and lower health-care costs in the USA, albeit that the costs of the intervention were not estimated (Cohen et al., 2006).Weekly group activity sessions led by occupational therapists in Canada significantly improved mental and physical health outcomes compared with participation in regular group social activities only. The incremental cost per QALY gained from a health and social care perspective was also considered to be costeffective (Hay et al., 2002). In Finland, a trial of psychosocial group therapy for older people identified to be lonely was also reported to be effective with a net mean reduction in healthcare costs per participant of E943 (Pitkala et al., 2009).
DISCUSSION
Our review indicates that there is an economic evidence base in all of the areas examined by DataPrev for some interventions to promote mental health and well-being in some very specific contexts and settings. In addition, we were able to identify published protocols of additional economic studies now underway. However, much of the existing economic literature that is available was beyond the scope of this review as it focused on actions targeted at the prevention of further deterioration, as well as the alleviation of problems in people already identified as having clinical threshold levels of mental disorder. This is consistent with the findings of previous reviews (Zechmeister et al., 2008).
One important limitation of our review was the restriction to English language only materials, although papers in other languages that had abstracts in English were included in the review. Certainly the overwhelming majority of material that we found came from English-speaking countries, but this is consistent with previous reviews of economic evaluations of public health interventions where no language restrictions were applied (McDaid and Needle, 2009).
We will have missed relevant studies concerning workplace interventions that have been published in diverse corporate literature with apparent positive returns on investment, but with insufficient information to be included in this review (Price Waterhouse Coopers, 2008). This includes case studies on the UK Health, Work and Wellbeing website looking at four large and small companies in the pharmaceutical, hotel and leisure, transport and manufacturing sectors. All report some positive impacts on absenteeism and/or staff retention rates. In the case of London Underground, for example, a return of 8:1 on investment in a stress management programme was reported (http://www.dwp.gov.uk/health-work-and-wellbeing/case-studies/).
Great caution must be exercised in drawing any firm conclusions on the economic case for investment, but the case for action in childhood or targeted at mothers appears strong. The economic consequences of poor mental health across different sectors and persisting into adulthood mean that effective health visiting and parenting programmes can have very favourable cost -benefit ratios; all economic analyses reported here from a societal perspective were cost-effective. Narrower perspectives adopted in some other child focused studies where evidence of effect was found, for instance from a health or education i134 D. McDaid and A-La Park perspective alone, may undervalue the potential case for action.
Nine of the ten economic analysis set in the workplace reported favourable outcomes. Most of these studies looked solely at the impacts for employers, either in terms of paying for the health care of their employees or dealing with absenteeism and poor performance at work. No studies looking solely at the benefits of organizational level actions to promote well-being and mental health were found. Given that there is a literature on the effectiveness of some of these measures, there is scope for modelling work to look at the potential economic costs and benefits of these measures. Of the 10 studies looking at programmes for older people, 3 were found to have little chance of being costeffective, but reasonable cost-effectiveness was reported for some group activities and home visiting activities.
In all areas we were able to identify published studies where no evidence of effect was found; these are also critical in helping to ensure resources are not used inappropriately. It is also the case that there has been little incentive to undertake formal economic evaluations of very low cost, but effective interventions, especially where costs are largely not borne by the public purse. One example of this are initiatives, often initially evaluated in low-and middle-income country contexts to promote skin-to-skin touch between mothers and their new born, where the principal cost is the time that the mother spends with her infant (Moore et al., 2007;Maulik and Darmstadt, 2009).
Going forward our analysis of the methodological quality of studies suggests much room for improvement. While high-quality analyses were identified, most studies failed to separate presentation of data on resources used to deliver interventions from the costs of these resources. Few studies undertook more than a very cursory sensitivity analysis to account for uncertainty around estimates of effect and cost. There was little discussion of the distributional impacts of interventions, an issue that is of particular relevance in the context of public health and health promotion interventions, where engagement and uptake can be critical to effectiveness (McDaid and Sassi, 2010).
There is also a need for more common and consistent endpoints to improve comparability across different interventions and country settings. Reliance solely on topic-specific outcomes such the cost per unit improvement of maternal sensitivity or a reduction in loneliness mean that it is difficult to compare the case for different potential areas of intervention. One key challenge in economic analysis going forward is to develop measures that can adequately capture the benefits of improved mental wellbeing. The principal quality of life measure reported in studies here, the QALY, was designed to identify the benefits of the absence of illness rather than well-being. Work on other approaches to well-being is underway; but in the meantime making use of validated well-being instruments such as the Warwick-Edinburgh Mental Wellbeing Scale (Tennant et al., 2007), alongside instruments used to value QALYs, such as the EQ-5D or SF-36, is merited. None of the cost -benefit analysis reported in this paper has elicited direct values for positive mental health: indeed the difficulty in putting a monetary value on well-being for cost -benefit analyses has been noted (Aos et al., 2004). Another issue is that despite the links between poor physical and poor mental health, little economic analysis has focused on the economic case for preventing co-morbidity, for instance on the prevention of depression to promote cardiovascular health. This is another area that economists might explore further.
More use can also be made of economic modelling in the short term to help strengthen the evidence base for investing in mental health and well-being. Such an approach has recently been used to help inform policy making on the case for prevention of various mental health problems in both England and Australia (Knapp et al., 2011;Mihalopoulos et al., 2011). The DataPrev project has demonstrated that there is a substantial evidence base on effective interventions; most of these have not been subject to economic evaluation. Working with programme implementers to determine resource requirements, costs of delivery and any necessary local adaptations, economic models could be used to determine the likelihood that interventions are likely to be cost-effective in different contexts, and over different time periods. They can also be used to look at the case for investing in multi-level approaches to promotion and prevention, with some interventions targeted at the general population and others targeted solely at high-risk groups. Published examples of this approach include the Triple P programme for children (Mihalopoulos et al., Investing in mental health and well-being i135 2007;Foster et al., 2008;Van't Veer-Tazelaar et al., 2010) and stepped care for older people. Such models could also factor in key critical factors such as probability of uptake and continued engagement by different population groups.
This work was supported by the European Commission
Sixth Framework Research Programme. Contract SP5A-CT-2007-044145. Funding to pay the Open Access publication charges for this article was provided by the LSE Institutional Publication Fund. | 9,468.2 | 2011-12-01T00:00:00.000 | [
"Economics",
"Medicine",
"Psychology"
] |
Search for time-reversal-invariance violation in double polarized antiproton-deuteron scattering
Apart from the $pd$ reaction also the scattering of antiprotons with transversal polarization $p_y^p$ on deuterons with tensor polarization $P_{xz}$ provides a null-test signal for time-reversal-invariance violating but parity conserving effects. Assuming that the time-reversal-invariance violating $\bar NN$ interaction contains the same operator structure as the $NN$ interaction, we discuss the energy dependence of the null-test signal in $\bar pd$ scattering on the basis of a calculation within the spin-dependent Glauber theory at beam energies of 50-300 MeV.
Introduction
Under CPT symmetry time-reversal-invariance violating but parity conserving (TVPC) forces are considered as a possible source of CP-invariance violation, which is required to account for the matterantimatter asymmetry in the universe [1]. In contrast to effects from time-reversal-invariance violation together with parity violation such as a permanent electric dipole moment (EDM) of elementary particles, so far much less attention was paid to TVPC effects. The reason why TVPC effects are interesting is that experimental limits on them are still rather weak, in particular, considerably weaker than those for the EDM.
Since the intensity of TVPC interactions within the standard model is extremely small [2], an observation of any effects at the present accuracy level of experiments would be a direct indication of physics beyond the standard model. Indeed a pertinent measurement is planned at the COSY accelerator in the Research Center in Jülich [3]. The observable in question is the integrated cross section for scattering of protons with transversal polarization p p y on deuterons with tensor polarization P xz . It provides a null-test signal for TVPC effects [4] and it will be measured in pd scattering at 135 MeV [3]. Theoretical studies of the energy dependence of the expected signal were performed at energies of the planned experiment [5][6][7][8][9][10][11][12] on the basis of the spin-dependent Glauber theory and demonstrate several unexpected effects. Among them are (i) the absense of the contribution from the lowest-mass meson-exchange (ρ meson) in the TVPC NN interaction, caused by its specific isospin, spin and momentum dependence; (ii) a strong impact of the deuteron D-wave on the null-test signal due to a destructive interference between the S -and D-wave contributions, even for zero transferred 3momentum; (iii) oscillating behaviour of the null-test signal as a function of the beam energy, i.e. the vanishing of the TVPC signal at some specific energies is possible even when the TVPC interaction itself is nonzero; (iv) a very small influence of the Coulomb interaction on the TVPC term of the pd forward scattering amplitude g. Furthermore, certain relations between differential observables of elastic pd scattering caused by time-reversal-invariance requirements were obtained and the degree of their violation by TVPC NN forces was studied [13,14].
Since the spin structure of the amplitude for pd-andpd elastic scattering is the same, it is obvious that the integrated cross section for scattering of a polarized (pp y ) antiproton on tensor polarized (P xz ) deuterons also provides a null-test signal for TVPC effects. Furthermore, the TVPCNN amplitude for elastic scattering contains the same operator structures as the one for TVPC NN elastic scattering, except for the charge-exchange terms. Therefore, the formalism developed in Refs. [7,8,11] within the Glauber theory for the calculation of the null-test signal in pd scattering can be straightforwardly applied topd scattering too. However, due to differences in the hadronic part of the pN andpN scattering amplitudes and also in the electromagnetic interactions, the energy dependence of the null test signal in pd andpd interaction has to be different. In the present work the energy dependence of the null-test signal inpd scattering is studied on the basis of calculations within the spin-dependent Glauber theory using the spin-dependentpN amplitudes from a recent partial wave analysis ofpp scattering [15].
Null-test signal for time-reversal-invariance violation
The total cross section forpd scattering with TVPC forces included can be written in the same form as for pd scattering [7] Here pp (p d ) is the vector polarization of the initial antiproton (deuteron), P zz and P xz are the tensor polarizations of the deuteron, and pp y is the transversal component of the antiproton vector polarization. The OZ axis is directed along the beam direction m, the OY axis is directed along the vector polarization of the antiproton beam pp and the OX axis is chosen to form a right-handed reference frame. The integrated cross sections σ t i (i = 0, 1, 2, 3) are those which arise from a standard timereversal invariant and parity conserving interaction, while the last term σ appears only in the presence of the TVPC interactions and constitutes the TVPC null-test signal. The result (1) can be derived using phenomenologicalpd forward scattering amplitudes and the generalized optical theorem.
The evaluation of the integrated cross sections σ t i and σ at beam energies > 100 MeV can be done on the basis of the spin-dependent Glauber theory ofpd scattering which is formulated similarly to the theory of pd scattering given in Ref. [16]. Indeed, as shown in Ref. [17], this theory allows one to describe rather well available data on differential spin observables of pd scattering in the forward hemisphere at beam energies of 135 − 200 MeV. For the antiproton-deuteron scattering this theory can be applied at even lower energies due to the presence of strong annihilation effects. In the Glauber theory one uses the elastic (on-shell)NN scattering amplitudes as input. Hadronic amplitudes of thē pN scattering are taken here in the same form as for pN scattering [16] M N (p, q; σ, σ N ) = A N + C N σn + C ′ N σ Nn + B N (σk)(σ Nk ) + whereq,k andn are defined as unit vectors along the vectors q = (p−p ′ ), k = (p+p ′ ) and n = [k×q], respectively; p (p ′ ) is the initial (final) antiproton momentum.
In general, the TVPC NN interaction contains 18 different terms [18]. In the case of the onshell NN scattering amplitude there are only three terms with different (independent) spin-momentum structures. In the present study we consider the following two terms for the TVPC (on-shell) t-matrix of elasticpN scattering which have the same structure as those in TVPC pN scattering Here σ (σ N ) is the Pauli matrix acting on the spin state of the antiproton (nucleon N = p, n) and τ (τ N ) is the isospin matrix acting on the isospin state of the antiproton (nucleon). The momenta q and k were already defined above in the context of Eq. (2). Both terms in Eq. (3), h N and g N , occur in the TVPC pn interaction. The TVPC pN scattering amplitude contains also the charge-exchange term which describes the elastic transitions pn → np and np → pn. Within a picture of one-mesonexchange interaction this g ′ -term corresponds to the charged ρ-meson exchange [19]. The same term (4) corresponds to the charge-exchange processespp →nn ornn →pp. However, in contrast to pn scattering these processes are inelastic and therefore the operation of time-reversal invariance transforms, for example, thepp →nn amplitude to thenn →pp amplitude and does not impose any restrictions on these amplitudes. The h N -term in Eq. (3) can be associated with the axial h 1 -meson exchange. As shown in Ref. [19], contributions of the π-and σ-meson to the TVPC NN interaction are excluded, which is obviously true for the TVPCNN interaction as well.
TVPC amplitude ofpd forward scattering
One can write thepd forward elastic scattering amplitude in general form taking into account the TVPCNN interactions, as it was done for pd elastic scattering [7,17], and then apply the generalized optical theorem to derive Eq. (1) for the totalpd scattering cross section. As in Ref. [7], the integrated cross section σ is related to the TVPC term g of thepd forward elastic scattering amplitude by σ = −4 √ π Im 2 3 g. Furthermore, the TVPC forward amplitude ofpd elastic scattering g can be found within the Glauber theory [7]. We consider the h N -and g N -terms and take into account both the Sand D-wave components of the deuteron. Taking into account that the g N -term is excluded in the processpn →pn due to the isospin operator in Eq. (3), we obtain the following result for the TVPC forward amplitude from the corresponding equation in Ref. [11]: Here S ( j) i are the elastic form factors of the deuteron defined in Ref. [11]. The first term in the (big) squared brackets in Eq. (5), S (0) 0 (q), corresponds to the S -wave approximation, the second term, S (1) 2 (q), accounts for the S -D interference, and the last three terms contain the pure D-wave contributions. As was shown in Ref. [11], the contribution of the g ′ -term to the null-test signal vanishes in pd scattering due to the specific spin-isospin structure of the g ′ -interaction. Formally, for the same reason the charge-exchange g ′ -term given by Eq. (4) vanishes in thepd forward elastic scattering amplitude.
In the first theoretical work [5] where the null-test signal was calculated within the impulse approximation, the Coulomb interaction was not considered. In Ref. [6] Faddeev calculations were performed, but only for nd scattering and at rather low energies of ∼ 100 keV. The Coulomb interaction was taken into account for the first time in Ref. [7] in a calculation of the null-test signal of pd scattering within Glauber theory and found to be negligible. A similar result was found in Ref. [20] using Faddeev calculations.
Numerical results
Results of numerical calculations of the energy dependence of the null test-signal for the h-term are presented in Fig. 1, in units of the unknown TVPC coupling strength. One can see from this figure that the deuteron S -wave contribution (dashed line) leads to a smooth energy dependence and has a node at an antiproton beam energy of about 50 MeV. The inclusion of the D-wave changes this behaviour considerably (solid line) due to a destructive S -D interference (cf. dash-dotted line). As a result, a second zero of the null-test signal σ appears at higher energies, i.e. at T ≈ 300 MeV. The maximal value of σ is expected at 100 − 150 MeV. Note that the actual position of the nodes changes only slightly when deuteron wave functions from other NN models are used for the calculation.
Let us consider possible spurious effects that could mimic a TVPC signal. One source for a spurious signal is associated with a nonzero deuteron vector polarization p y d 0 (in the direction of the incident-proton-beam polarization p p ). In this case, the term σ 1 Pp y p d y in Eq. (1) contributes to the 4 EPJ Web of Conferences 181, 01015 (2018) https://doi.org/10.1051/epjconf/201818101015 EXA2017 asymmetry corresponding to the difference of the event counting rates for the cases of pp y P xz > 0 and p p y P xz < 0 (with the fixed sign of P xz ), which is planned to be measured at COSY [3]. According to our calculations, the integrated cross section σ 1 could be equal to zero at antiproton beam energies of ∼ 100 MeV (see results for the JülichNN interaction model in Refs. [21,22]). Therefore, at this energy the spurios signal caused by a nonzero value of the deuteron vector polariziation p d y could be minimized.
Concluding remarks
We have performed a study of time-reversal-invariance violating but parity conserving effects in antiproton-deuteron scattering. Specifically, we have evaluated the null-test TVPC signal for scattering of antiprotons with transversal polarization p p y on deuterons with tensor polarization P xz on the basis of the spin-dependent Glauber theory. The observed effects turned out to be similar to those in pd scattering: (i) There is a strong impact of the deuteron D-wave on the null-test signal that arises from a destructive interference between the S -and D-wave contributions; (ii) There is an oscillating behaviour of the null-test signal as a function of the beam energy. Accordingly, it is possible that the signal for TVPC effects is zero at some specific energies, even when the TVPC interaction itself is nonzero. | 2,860.2 | 2017-12-14T00:00:00.000 | [
"Physics"
] |
Improving Optical-Wireless CDMA System Performance in Industrial Environment with Timing Jitters
In this paper, a novel analytical model is proposed and formulated to quantify the timing jitters, introduced by environmental changes, in optical-wireless code-division multiple access systems. The model divides every chip in an optical codeword into multiple equal intervals, and each pulse in the codeword can be randomly shifted to one of these sub-chip positions in order to account for the effect of the timing jitters. Our study shows that the new model can make a good use of the time skew of pulses in optical codewords and unconventionally improve O-CDMA performance under a certain condition.
Introduction
Direct-detection (or so-called incoherent) optical code division multiple access (O-CDMA) has been studied for applications in fiber-optic and optical-wireless multipleaccess systems and networks because of its desirable features, such as flexible bandwidth utilization, asynchronous access without the need of precise coordination, efficiency in bursty traffic, and dynamic optical-channel sharing without complex scheduling (1)(2)(3)(4) .While CDMA has already been used in wireless communications, its usefulness in industrial/manufacturing plants that see strong electromagnetic interference (EMI) is moot.As optical technology is immune to EMI, the use of CDMA in optical-wireless computer/communications/control networks in such strong-EMI environment becomes attractive.Also, optical-wireless CDMA can also provide mobility and ease of set-up/tear- Fig. 1 Example of an optical-wireless CDMA system using free space as the multi-access optical channel, suitable for manufacturing plants that have strong EMI and prefer node mobility.
down (5)(6)(7)(8) .Figure 1 shows an example of an optical-wireless CDMA system model, in which optical codeword of each user (or node) is transmitted via free space for mobility and EMI immunity.Assume that every user sends its data bits in on off keying (OOK) modulation format.In the transmitter of a node, a gated optical pulse, representing the transmission of a data bit of 1, is first encoded into the address codeword of the receiver of the intended node.The structure of the 1-D/2-D optical encoders depends on the 1-D or 2-D coding scheme in use.The codeword is then transmitted onto the diffuser(s) on the ceiling and, in turn, distributed to all nodes.In each receiver, the optical decoders serve as inverted filters of the optical encoders.Each decoder matches the time positions (and wavelengths, if 2-D wavelength-time codes are used) of the pulses of arriving codewords with its address signature.A hard-limiter can be placed at the front end of each optical decoder to reduce MAI-localization and near-far problems (9,10) .
Two important issues in optical-wireless CDMA are on the codeword arrival-time tracking and integrity of the op-tical pulses (i.e., precisely sitting in their associated chips or time slots) within each codeword.It is commonly assumed that the arrival time of the codewords from a node is fixed after the timing of the node has been established in the system.Nevertheless, physical or environmental changes, such as temperature fluctuations, may cause timing jitters and, in turn, generate a slight mismatch in the codeword's arrival time at the intended receiver and also create time skew on the optical pulses within the codeword.For example, the effect of environmental temperature fluctuations to the performance of a long-haul fiber-optic CDMA system was studied by Osadola, et al. in (11) .They included fiber thermal coefficient in the analysis and modeled the effect of temperature variations as time skew on the optical pulses in codewords.The distortion of autocorrelation peaks and, in turn, worsening of system performance were formulated as a function of the amount of time skew introduced by temperature fluctuations and distance traveled.Even though their study was not directly applicable, it raised an important fact that any time skew on the optical pulses within transmitted codewords would cause performance degradation in optical-wireless CDMA systems as well.
However, their analytical model did not take into account that the cross-correlation properties of the optical codewords in use could also be changed by time skew, in addition to the autocorrelation peaks.Our study in Section 2 shows that time skew indeed worsens the cross-correlation values of the optical codes in use.In Section 3, a novel analytical model, which can be used to complement the model in (11) , is formulated.Our results imply that time skew of optical pulses, caused by environmental changes or timing jitters, can constructively be exploited to improve system performance if there exists a feedback mechanism in the receiver to reconstruct the original autocorrelation peak.This new finding is unconventional in the sense that an O-CDMA system can make good use of this deleterious effect to improve performance, rather than harming it, opposite to the finding in (11) .Finally, the new model is validated by numerical examples and computer simulation in Section 4.
Quantifying Time Skew and Its Effect to Cross Correlations
In incoherent O-CDMA with OOK modulation, each user conveys the address codeword of its intended receiver whenever a data bit of 1 is transmitted, but nothing is conveyed for bits of 0. Assume the use of a family of (L × N, w, (2, 4) .In general, every 2-D codeword, say codeword i, in the code set can be represented its w pulses in form of w ordered pairs such that C i = [(λ 0 , t 0 ), (λ 1 , t 1 ), …, (λ w-1 , t w-1 )], where each ordered pair denotes that the pulse of wavelength λ j is located in time-slot (or chip) position t j ∈{0, 1, …, N -1} for all j∈{0, 1, …, w -1}.For the case of 1-D optical codes, L = 1 and all w pulses in every codeword in the code set use one identical wavelength such that λ 0 = λ 1 = …=λ w-1 .
Due to temperature fluctuations, the time skew of optical pulses can here be quantified as random time shift.In our model, the chips (or time slots) of the optical pulses in every codeword is subdivided into s sub-chips of equal width, where s > 1 is an integer.Each of these w pulses can be randomly shifted to start at any one of these s sub-chips from its original chip.Let the width of every chip be equal to 1.Then, the width of every sub-chip is 1/s.Also let the time delay of the jth pulse (for 1-D codes) or the pulse of the jth wavelength (for 2-D codes) created by an independent random sub-chip shift be denoted as τ j ∈{0, 1/s, 2/s, …, (s-1)/s} for all j∈[0, w-1].So, the ordered pairs of the "shifted" copy of C i become [(λ 0 , t 0 +τ 0 ), (λ 1 , t 1 +τ 1 ), …, (λ w-1 , t w-1 +τ w-1 )].
In the above example, the cross-correlation process between the shifted copies of two codewords can created as large as 2 pulse-overlaps (or so-called hits), even though the original cross-correlation function of the codes was at most 1 (i.e., λ c =1).In general, for a given s>1, it is found that the maximum cross-correlation value of the (w ×p 1 p 2 … p k , w, 1) CHPCs is upper bounded by where w ≥ 3.This upper bound can also be applied to any 1-D and 2-D optical codes that each of the codewords uses at most one pulse per wavelength and per chip, and no wavelengths are used more than once within each codeword.
Performance Analysis
With s >1, the optical pulses in each codeword are assumed to be independently and randomly shifted by some integral multiples of 1/s of a chip, the involved crosscorrelation function now carries values between the discrete cross-correlation function in the "pure" chip-synchronous case and the continuous cross-correlation function in the "pure" chip-asynchronous case (2,4) .The new cross-correlation function is still discrete but now with fractional values taking from the set of {0, 1/s, 2/s, …, 3-2/s}, according to (1).Taking these fractional values into consideration, the new hit probabilities can here be formulated as which are indexed by kʹ and lʹ∈{0, 1/s, 2/s, …, 3-2/s}.The factor 1/2 is due to OOK, ϕ -1 represents the possible number of interfering codewords, out of a total of ϕ codewords in the code set, N represents the number of possible time shifts in a codeword of length N, and s 2w represents all the possible pulse matching situations between two correlating codewords as there exist a total of s w pulse-shift positions in a codeword of weigh w.The term h k',l' denotes the number of times of getting a kʹ-hit in the preceding sub-chip and a lʹ-hit in the current sub-chip, where kʹ and lʹ∈{0, 1/s, 2/s, …, 3-2/s}, which depends on the optical codes in use and can be computed numerically.
For example, if the (w × p 1 p 2 …p m , w, 1) CHPCs with s= 2 are used, ϕ=p 1 p 2 …p m for a given integer m ≥ 1 and p m ≥ p m-1 ≥ …≥ p 2 ≥ p 1 ≥ w.The term h k',l' , for kʹ and lʹ∈{0, 1/2, 1, 3/2, 2}, can be computed by i) first randomly pick two correlating codewords from the code set; ii) then build the s w possible sub-chip shifts in the w pulses of these two codewords; iii) additionally build the N possible cyclic chip shifts in one of the codeword; iv) for every combination of the s w sub-chip and N chip shifts, count the number of times of getting a kʹ-hit in the preceding sub-chip and a lʹ-hit in the present sub-chip in the cross-correlation function, and then add to the associated h k',l' term; v) finally repeat the above steps with all possible combinations of two correlating codewords in the code set.
In general, the chip-synchronous error probability P e,syn,s>1 caused by time skew with s >1 can be derived as (14) where K denote the number of simultaneous users.The rational of the derivation of (3) follows that of the pure chip-asynchronous analysis (2,4) .Under the pure chipasynchronous assumption, the cross-correlation value becomes a continuous function of time because the value is the amount of (partial) pulse overlap due to the relative time-shift between the two correlating codewords (4, p. 36) .As a result, the "asynchronous" cross-correlation function involves partial overlap of pulses found in two consecutive chips.Thus, q k,l is defined as the probability of the crosscorrelation value in the preceding chip equal to k∈[0, λ c =1] and the crosscorrelation value in the present chip equal to l ∈[0, λ c =1], as if it is under the pure chip-synchronous assumption.For λ c =1 optical codes, these hit probabilities are related by q 1,0 =q 0,1 , q 1 =w 2 /(2LN), q 1,1 =w(w-1)/[2N(N-1)], q 1 =0.5Σ i=0 Σ j=0 (i+j)q i,j, and Σ i=0 Σ j=0 q i,j =1 (2,4) .
Numerical Results
In Figure 3, the hard-limiting error probabilities, P e,syn,s>1 of (3), P e,asyn of (5), and P e,syn of (6), of the (L×N, w, 1) CHPCs are plotted against the number of simultaneous users K, where w=L={3, 5, 7}, N={9, 25, 49}, and s={2, 3, 4}.In general, the performance (i.e., P e ) gets worse as K increases due to stronger MAI.The code performance improves with L, N, w, and s because the increment of L or N reduces the hit probabilities, the increment of w increases the autocorrelation peak, and the hit probabilities reduced with increasing s.For a given set of (w, N) values, the dotted curves of P e,syn,s>1 are bounded by the solid curve of P e,syn and the dashed curve of P e,asyn .This is because these two curves correspond to the extreme cases of pure chipsynchronism and chip-asynchronism, respectively.Also shown in Figure 3 are the computer-simulation results (i.e., asterisks) of the hard-limiting error probabilities with the same code parameters as the corresponding theoretical (solid, dotted, and dashed) curves.Both theoretical and computer-simulation results (i.e., dotted curves vs. asterisks) match closely at various s and (w, N) values, thus validating the accuracy of the analytical model derive in (3). Figure 4 plots the hard-limiting error probability, P e,syn,s>1 of (3), of the (5×25, 5, 1) CHPCs as a function of s with K=13.In general, the red curve achieves better performance (i.e., lower P e ) as s increases because more sub-chips increases the number of possible locations for the pulses in each codeword, thus reducing the hit probabilities.Also shown in the figure are the computer-simulation results under the chip-synchronous (represented as circles) and chip-asynchronous (represented as asterisks) assumptions.The circles match closely with the red curve of P e,syn,s>1 from (3).The asterisks represents the results when chipasynchronism is applied in the simulation, in which a user begins codeword transmission at any time.In summary, by using different number of sub-chips (i.e., s∈ [2, 16]), the error probability is improved by 2-3 orders of magnitude as shown in the red curve.From the chip-asynchronous simulation results (i.e., asterisks), the error probability of s= 16 is about 6 times better than that of s=1.
Conclusions
In this paper, a new analytical model for the time skew of optical pulses in O-CDMA codewords due to environmental changes was investigated.The hard-limiting performance of such a O-CDMA system was formulated, illustrated with a numerical example, and validated by computer simulation.Our study showed that larger number of sub-chips, s, increased the number of possible locations for the pulses in a codeword, thus reducing the hit probabilities and amount of MAI contributed by other simultaneous users.As a result, the performance improved as s got larger.This new finding is unconventional in the sense that an O-CDMA system can make use of the time skew to improve performance. | 3,168.2 | 2015-02-05T00:00:00.000 | [
"Engineering",
"Physics",
"Computer Science"
] |
Mechanical Properties and Structures of Clay-Polyelectrolyte Blend Hydrogels
Our recent studies have shown that the hydrogels prepared by blending clay, a dispersant of clay, and a polyelectrolyte (sodium polyacrylate (PAAS)) possess excellent mechanical properties. In order to clarify the mechanism of the toughness, we have so far investigated the effects of the composition, molecular mass of the polymer, and kinds of polymers on the mechanical properties. This study has focused upon the mechanical properties and structures of the clay/PAAS gels using three kinds of smectite clay minerals such as synthetic hectorite (laponite XLG), saponite (sumecton-SA), montmorillonite (kunipia-F), whose particle size becomes larger according to the sequence. Laponite/PAAS and sumecton/PAAS gels were quite tough for high compression, whereas kunipia-F/PAAS did not gelate. In comparison between sumecton/PAAS gel and laponite/PAAS gel, the mechanical property of the former gel was poorer than that of the latter gel due to the inhomogeneous distribution of clay platelets in the gel. Synchrotron small-angle X-ray scattering experiments revealed that their clay platelets laid down in the stretching direction under elongation. Furthermore, it was found that sumecton/PAAS gel under elongation was arranged with an interparticle distance of ~6.3 nm in the direction perpendicular to the stretching. Such local ordering under elongation may originate in local aggregation of sumecton platelets in the original state without elongation.
Introduction
Hydrogels composed of clay and polymer have attracted widespread interest due to their potential applications in various fields. Incorporation of the clay into polymer hydrogel results in improvement of the mechanical properties [1][2][3]. Most of such clay-polymer nanocomposite hydrogels with excellent mechanical performance have been prepared by the in situ polymerization method [4][5][6][7][8]. Their studies have shown that the tensile stress increased with the increase of clay concentration [9,10] and interactions between clay and polymer are attributed to noncovalent bonds, e.g., hydrogen bonds between the amide groups and silanol groups (Si-OH) and/or siloxane units (Si-O-Si) play an important role on the cross-link of clay-poly(alkyl acrylamide) nanocomposite hydrogel [11]. In addition, the effect of clay type (synthetic hectorite, fluorinated hectorite, natural montmorillonite, or sepiolite magnesium silicate) on the mechanical performance of the nanocomposite hydrogels has been examined, so that it has been shown that synthetic hectorite (laponite XLG)-the polymer nanocomposite hydrogel has the best mechanical performance [11].
Recently we have succeeded in fabricating mechanically tough hydrogels composed of clay, a dispersant of clay, and an ultrahigh molecular mass polyelectrolyte such as sodium polyacrylate (PAAS), which have been prepared by blending them [12]. The clay/PAAS hydrogels are tough and Gels 2018, 4, 71 2 of 9 thus suitable for high compression, and have a very large swelling ratio. So far we investigated effects of molecular mass of polymers and the composition on the mechanical properties of the hydrogels. As a consequence, it has been clarified that the factors such as dispersion of clay platelets, use of ultrahigh molecular mass polymers higher than a few million, and favorable interactions between clay and polymers are important for accomplishment of the toughness [13][14][15].
In this study, we investigated the mechanical properties and structures of the clay/PAAS blend hydrogels using three kinds of layered silicate minerals belonging to smectite group, such as synthetic hectorite, synthetic saponite, and montmorillonite. Smectite clay minerals have a 2:1 layer structure, in which an octahedral sheet is sandwiched between two tetrahedral sheets. Although layers of these clay minerals have nearly the same thickness of~1 nm, the dimensions are very different. As mentioned above, although the effect of the clay type on the mechanical properties of the clay-polymer nanocomposite hydrogels synthesized by the in situ polymerization method has been investigated, no studies have made for the nanocomposite hydrogels prepared by blending (blend hydrogels). In this study, our peculiar attention is to investigate the effect of clay size on the mechanical properties of the clay/PAAS blend hydrogels.
Results and Discussion
Firstly, we estimated the size of the clay platelets with synchrotron small-angle X-ray scattering (SAXS) for dilute clay aqueous dispersions. Assuming that the interference effects between clay platelets are negligible for a very dilute clay dispersion, the scattering intensity I(q) is expressed by where n clay and ∆ρ are the number density of a clay platelet with a volume of V clay and the scattering contrast factor, respectively. Here P(q) represents the form factor of randomly distributed cylindrical particles with the radius R and thickness 2H [16,17] P(q) = 4 π/2 0 sin 2 (qH cos β) (qH cos β) 2 where J 1 is the Bessel function of the first order, and β represents the angle between q and the axis of the disk. Figure 1 depicts SAXS curves for dilute clay aqueous solutions, (a) a 0.3 wt % laponite solution and (b) a 0.3 wt % sumecton solution. The curves fitted with Equation (2) are in good agreement with the scattering data. The fitting analysis showed that the size of sumecton (R = 23 nm) was larger than that of laponite (R = 14 nm).
Gels 2018, 4, x FOR PEER REVIEW 2 of 9 (PAAS), which have been prepared by blending them [12]. The clay/PAAS hydrogels are tough and thus suitable for high compression, and have a very large swelling ratio. So far we investigated effects of molecular mass of polymers and the composition on the mechanical properties of the hydrogels. As a consequence, it has been clarified that the factors such as dispersion of clay platelets, use of ultrahigh molecular mass polymers higher than a few million, and favorable interactions between clay and polymers are important for accomplishment of the toughness [13][14][15].
In this study, we investigated the mechanical properties and structures of the clay/PAAS blend hydrogels using three kinds of layered silicate minerals belonging to smectite group, such as synthetic hectorite, synthetic saponite, and montmorillonite. Smectite clay minerals have a 2:1 layer structure, in which an octahedral sheet is sandwiched between two tetrahedral sheets. Although layers of these clay minerals have nearly the same thickness of ~1 nm, the dimensions are very different. As mentioned above, although the effect of the clay type on the mechanical properties of the clay-polymer nanocomposite hydrogels synthesized by the in situ polymerization method has been investigated, no studies have made for the nanocomposite hydrogels prepared by blending (blend hydrogels). In this study, our peculiar attention is to investigate the effect of clay size on the mechanical properties of the clay/PAAS blend hydrogels.
Results and Discussion
Firstly, we estimated the size of the clay platelets with synchrotron small-angle X-ray scattering (SAXS) for dilute clay aqueous dispersions. Assuming that the interference effects between clay platelets are negligible for a very dilute clay dispersion, the scattering intensity I(q) is expressed by where nclay and ∆ρ are the number density of a clay platelet with a volume of Vclay and the scattering contrast factor, respectively. Here P(q) represents the form factor of randomly distributed cylindrical particles with the radius R and thickness 2H [16,17] (2) where J1 is the Bessel function of the first order, and β represents the angle between q and the axis of the disk. Next we investigated the mechanical properties of the clay/PAAS blend hydrogels. Figure 2 shows pictures of a 10 wt % sumecton/PAAS gel before, during, and after compression. The sumecton/PAAS gel almost recovered its initial shape after load removal, as shown in the pictures. The behavior is the same as the laponite/PAAS gel, whose pictures were shown in the previous study [12]. The elastic behavior of the sumecton/PAAS and laponite/PAAS gels suggests incorporation of clay platelets into the polymer matrix, i.e., formation of clay/polymer nanocomposite hydrogel, where clay platelets act as a multiple cross-linker. In practice, as the gels show an alkaline nature (pH ≈ 9.9), as shown in Table 1, PAAS exists as carboxylate ions in the gel, which is expected to adsorb upon the positively charged edge of clay platelets. The comparison of compressive properties of the sumecton/PAAS and laponite/PAAS gels showed that the compressive stress of the former gel was lower than that of the latter gel ( Figure 3). On the other hand, the mixture of PAAS and kunipia F with larger particle size (300~500 nm) [18] did not gelate. This result is different from that of the kunipia F-poly(N,N-dimethyl acrylamide) (PDMAA) gel prepared by the in situ polymerization method. Although the mechanical performance for the kunipia F-PDMAA nanocomposite hydrogel was lower than that of the laponite/PDMAA nanocomposite hydrogel, the former gel also showed high elongation [11]. Such a difference may come from the different preparation processes of the synthesized gel and the blend gel, i.e., as discussed in the previous paper [14], in the case of the synthesized gel, even if monomers before polymerization are added in the kunipia F dispersion, the resulting solution may not become inhomogeneous. Though the progress of polymerization may induce phase separation, or an inhomogeneous structure, the resulting gelation may give rise to freezing of the structure and prevent the inhomogeneous structure from forming. Contrary to this, in the case of the blend hydrogel, the mixture of kunipia F with large size and ultrahigh molecular mass PAAS has a tendency to become thermodynamically immiscible, because the combinatorial entropy of mixing is very small. Consequently, the effect of the clay size on the mechanical performance becomes more sensitive for the blend hydrogel. Thus, as the clay size is larger, the mechanical properties became poorer. The elastic moduli for laponite/PAAS and sumecton/PAAS gels are summarized in Table 1. The modulus of the sumecton/PAAS gel was slightly lower than that of the laponite/PAAS gel at 5 wt % clay concentration, whereas the former was much lower than the latter at 10 wt % clay concentration. The earlier studies have shown that an increase in clay concentration gives rise to a considerable increase in the elastic modulus of the gel, unless dispersion of clay platelets in the gel is significantly affected [1,12,14]. Therefore, the elastic modulus of the laponite/PAAS gel remarkably increases with the increase of the clay concentration. This behavior is attributed to the increase in the functionality of the gel that represents the number of the cross-linking points per clay platelet [12,14]. On the other hand, in the case of the sumecton/PAAS gel, the increase in clay concentration gave rise to only a slight increase in the elastic modulus. The result was attributed to poor dispersion of sumecton platelets at the higher clay concentration, as suggested from the results of transmittance in Table 1, which tends to reduce the mechanical performance. On the other hand, dispersion of sumecton platelets at the lower clay concentration may be moderately poor. As a matter of fact, the 5 wt % laponite/PAAS gel was transparent, whereas the 5 wt % sumecton/PAAS gel was translucent, as shown in Figure 4. Next, we examined the tensile properties of laponite/PAAS and sumecton/PAAS gels. Figure 5 depicts representative tensile stress-strain curves for the 5 wt % laponite/PAAS and 5 wt % sumecton/PAAS gels. The tensile strength and extension ratio for the laponite/PAAS gel was better than those of the sumecton/PAAS gel. In order to investigate structure of the gels during stretching, synchrotron SAXS experiments were performed. Although the clay/PAAS blend hydrogel is a four-component system and therefore the scattering intensity of the gel is given by the sum of the partial scattering functions between the respective components, in fact the partial scattering function for the clay-clay components is dominant in the X-ray scattering, as mentioned in the previous paper [14]. Namely, we can clearly see the structure of the clay platelets from the SAXS intensity of the gel. This is because the clay platelets are composed of heavier atoms, which have larger X-ray scattering length [17]. Figure 6 depicts two-dimensional patterns for laponite/PAAS and sumecton/PAAS gels during elongation. The SAXS patterns for both gels are isotropic before stretching (see the SAXS pattern at the extension Gels 2018, 4, 71 5 of 9 ratio λ = 1), reflecting that the structure is isotropic before stretching, whereas they show elliptic patterns with a longer axis in the vertical direction during stretching. The anisotropic character increased with the increase of the stretching ratio. The elliptic patterns suggest that the clay platelets lie down towards the stretching direction under elongation, considering an inverse relationship between the object size and the scattering angle in scattering theory [19]. As a matter of fact, the scattering intensity for both gels in the perpendicular direction was stronger than that in the parallel direction, as shown in Figures 7 and 8. These results also support the above interpretation. The scattering curves of the sumecton/PAAS gel in the perpendicular direction at high elongation had multiple small peaks or shoulders (Figure 9). The ratio of the peak position was 1:2:3 and the first peak was observed at q ≈ 0.1 Å −1 . Thus, sumecton/PAAS gel may be arranged in the perpendicular direction with an interparticle distance of~6.3 nm at high elongations. Such regular structure may reflect local aggregation of clay platelets in the original state (undeformed state), which contributes to lower the mechanical performance of the gel.
Sumecton (x = 10) 9.92 4.7 ± 0.3 1.10 Laponite (x = 10) 9.86 26 ± 7.3 60.7 In order to investigate structure of the gels during stretching, synchrotron SAXS experiments were performed. Although the clay/PAAS blend hydrogel is a four-component system and therefore the scattering intensity of the gel is given by the sum of the partial scattering functions between the respective components, in fact the partial scattering function for the clay-clay components is dominant in the X-ray scattering, as mentioned in the previous paper [14]. Namely, we can clearly see the structure of the clay platelets from the SAXS intensity of the gel. This is because the clay platelets are composed of heavier atoms, which have larger X-ray scattering length [17]. Figure 6 depicts two-dimensional patterns for laponite/PAAS and sumecton/PAAS gels during elongation. The SAXS patterns for both gels are isotropic before stretching (see the SAXS pattern at the extension
Conclusions
In this study, we examined the mechanical properties and structures of hydrogels composed of sumetite clay minerals and PAAS. The dimensions of the clay platelets significantly affects the mechanical performance of the gels, i.e., as the dimensions of the clay platelets becomes larger, the mechanical performances of the gels are lowered. The mechanical strength of sumecton/PAAS gel was slightly worse than that of laponite/PAAS gel. The lowering of the mechanical strength for the sumecton/PAAS hydrogel was found to come from the structural inhomogeneity in the hydrogel. The SAXS experiments under elongation revealed that the clay platelets for both laponite/PAAS and sumecton/PAAS gels laid down in the stretching direction. The latter gel showed a regular structure with an interparticle distance of ~6.3 nm in the direction perpendicular to the stretching, which may be originated from the local aggregation in the original state.
Sample and Sample Preparation
In this study, we used three kinds of laponite XLG (RockWood Ltd., Newry, UK), sumecton-SA, and Kunipia-F (Kunimine Industries Co., Ltd., Tokyo, Japan) as clay samples, a dispersant of clay platelets (tetrasodium pyrophosphate (TSPP)) and sodium polyacrylates (PAAS) with a weightaveraged molecular mass of 3.5 × 10 6 from Wako Pure Chemical Industries, Ltd., Osaka, Japan. The hydrogel was prepared in the same manner as described in the previous study [13,14], except for using a sonicator (QSONICA Q55, Waken Tech co., Ltd., Kyoto, Japan) for the dispersion of clay platelets. The final concentration of PAAS and TSPP in the hydrogels was 1 wt % and 0.5 wt %, respectively.
Transmittance and pH Measurements
Transmittance measurements were performed for laponite and sumecton gels with a thickness of 6 mm using a He-Ne laser with a wavelength of 632.8 nm. pH values of the gel surface were measured using a pH meter (Horiba, LAQUA F-71, Kyoto, Japan) with a pH electrode (Horiba, ISFET 0040-10D).
Compression and Tensile Measurements
Compression tests were conducted for the gels with a cylindrical shape of a diameter of 14 mm and a thickness of 8 mm using AIKOH Engineering 1305NR (Osaka, Japan) with a load cell of 20 N
Conclusions
In this study, we examined the mechanical properties and structures of hydrogels composed of sumetite clay minerals and PAAS. The dimensions of the clay platelets significantly affects the mechanical performance of the gels, i.e., as the dimensions of the clay platelets becomes larger, the mechanical performances of the gels are lowered. The mechanical strength of sumecton/PAAS gel was slightly worse than that of laponite/PAAS gel. The lowering of the mechanical strength for the sumecton/PAAS hydrogel was found to come from the structural inhomogeneity in the hydrogel. The SAXS experiments under elongation revealed that the clay platelets for both laponite/PAAS and sumecton/PAAS gels laid down in the stretching direction. The latter gel showed a regular structure with an interparticle distance of~6.3 nm in the direction perpendicular to the stretching, which may be originated from the local aggregation in the original state.
Sample and Sample Preparation
In this study, we used three kinds of laponite XLG (RockWood Ltd., Newry, UK), sumecton-SA, and Kunipia-F (Kunimine Industries Co., Ltd., Tokyo, Japan) as clay samples, a dispersant of clay platelets (tetrasodium pyrophosphate (TSPP)) and sodium polyacrylates (PAAS) with a weight-averaged molecular mass of 3.5 × 10 6 from Wako Pure Chemical Industries, Ltd., Osaka, Japan. The hydrogel was prepared in the same manner as described in the previous study [13,14], except for using a sonicator (QSONICA Q55, Waken Tech co., Ltd., Kyoto, Japan) for the dispersion of clay platelets. The final concentration of PAAS and TSPP in the hydrogels was 1 wt % and 0.5 wt %, respectively.
Transmittance and pH Measurements
Transmittance measurements were performed for laponite and sumecton gels with a thickness of 6 mm using a He-Ne laser with a wavelength of 632.8 nm. pH values of the gel surface were measured using a pH meter (Horiba, LAQUA F-71, Kyoto, Japan) with a pH electrode (Horiba, ISFET 0040-10D).
Compression and Tensile Measurements
Compression tests were conducted for the gels with a cylindrical shape of a diameter of 14 mm and a thickness of 8 mm using AIKOH Engineering 1305NR (Osaka, Japan) with a load cell of 20 N and FA1015B Force Analyzer Explorer III at a compression speed of 10 mm/min. Tensile measurements were performed with an ORIENTEC TENSILE TESTER STM-20 (Tokyo, Japan) at a compression speed of 10 mm/min. Stress was calculated using a cross-sectional area of an undeformed gel. The elastic modulus E was obtained from the slope of the stress-strain curves at the small strains, and the average value of E was evaluated from three tests.
Synchrotron Small-Angle X-ray Scattering (SAXS)
Synchrotron small-angle X-ray scattering (SAXS) was performed at beam line 6A and 10C of the photon factory (PF) of the High Energy Accelerator Research Organization (KEK) in Japan. The SAXS experiments were conducted with a wavelength of λ = 1.5 Å, and the samples for the measurements were put in an aluminum spacer with a thickness of 1 mm using a very thin Kapton film as a window. The SAXS data was collected by a two dimensional detector (PILATUS-2M or PILATUS-1M) and then circularly averaged to obtain the scattering curves as a function of q defined as follows, where θ denotes the scattering angle. The scattered intensity obtained thus was corrected for the background scattering, and then reduced to the absolute unit using a glassy carbon [20].
Author Contributions: H.T. analyzed the experimental data and wrote the paper. S.N. performed the experiments and analyzed the data.
Funding: This research was funded by JSPS KAKENHI Grant Number JP15K05242. | 4,544.4 | 2018-08-30T00:00:00.000 | [
"Materials Science"
] |
A Rare Presentation of Tuberculosis-Related Septic Shock
Septic shock with multi-organ dysfunction is an exceedingly rare, but known complication of untreated Mycobacterium tuberculosis (TB) infection. TB-associated cases of septic shock are predominantly reported in immunocompromised patients; however, it can manifest in a healthy individual if the infection is not treated. Through the interaction of lipoarabinomannan (LAM) on the mycobacterium cell wall with antigen-presenting cells, the bacteria may be able to survive in host cells for long periods of time. Without prompt treatment, TB may cause bronchiectasis and multi-organ failure. We report a case of a 24-year-old woman with untreated TB who developed widespread bronchiectasis and septic shock.
Introduction
There are over 10 million newly diagnosed tuberculosis (TB) cases worldwide. More specifically, the incidence of TB in Guyana is currently greater than 100 per 100,000 [1]. TB-related cases of septic shock are remarkably rare and mainly seen in immunocompromised hosts. The following is a case of a malnourished female from Guyana who presented with severe malnutrition and septic shock in the setting of untreated Mycobacterium tuberculosis.
Case Presentation
A 24-year-old Guyanese woman with no past medical history presented to the emergency department with shortness of breath and weight loss over the past three months. She had also noted six weeks of a productive cough, hemoptysis, abdominal pain, and night sweats. The patient did not seek medical attention or undergo any previous treatment and remained home despite symptoms. Minimal history was obtained through a family member; the patient had not travelled recently, was not currently employed, and had a history of previously diagnosed TB which was unclear.
In the emergency department, she was tachycardic to 140 bpm, hypotensive to 80/40 mmHg, tachypnic at 40 breaths per minute, and hypoxic to 82% on room air requiring a high-flow nasal cannula. She appeared ill and cachectic with a BMI of 10 kg/m2. The exam was further remarkable for accessory muscle use and diffuse coarse breath sounds in all lung fields. An ECG showed sinus tachycardia, lactate of 4 mmol/L, haemoglobin of 6.6 g/dL, and albumin of 2 g/dL. A chest X-ray was significant for diffuse cystic changes ( Figure 1).
FIGURE 2: Varicose bronchiectasis
CT scans of the chest with contrast showed extensive severe cystic and varicose bronchiectasis throughout the lung.
A right internal jugular central line was placed and she was given bolus fluids with minimal improvement in vital signs. She was started on broad-spectrum antibiotics and was admitted to the medical intensive care unit.
In the ICU, standard blood cultures remained negative. Further workup including HIV testing and a respiratory viral panel was negative. Given the high suspicion of TB-related septic shock, the patient was started on empiric parenteral rifampin, isoniazid, pyrazinamide, and ethambutol (RIPE) therapy on hospital day 1 of 2 while waiting for acid-fast sputum culture results, which was positive the following day. She was also started on glucocorticoids due to TB-related adrenal insufficiency. A bedside echo was performed in the ICU which demonstrated collapse of the left ventricular walls at end-systole with a decreased left ventricular end-diastolic pressure, reduced cardiac output, a significantly collapsed inferior vena cava, and no signs of pericardial effusion. Overnight, the patient became severely hypoxic requiring intubation. Hours later she developed refractory hypotension despite four different vasopressors. A new chest X-ray redemonstrated bronchiectasis with extensive consolidations bilaterally ( Figure 3).
FIGURE 3: ICU chest X-ray with extensive cavitation and consolidation
A chest X-ray was taken in the ICU redemonstrating cystic bronchiectasis and severe cavitary disease with extensive consolidations bilaterally.
She became pulseless shortly afterwards and advanced cardiac life support (ACLS) protocol was initiated. The patient remained in asystole throughout the arrest. Unfortunately, we were unable to achieve return of spontaneous circulation (ROSC) and the patient expired.
Discussion
TB is a disease caused by Mycobacterium tuberculosis which is a unique bacterium with mycolic acid on the cell surface making it difficult to gram stain [2]. In a retrospective global cohort study, Kethireddy et al. were able to isolate TB as the causative agent of septic shock in only 1% of cases [3]. Most patients can recover from mild sepsis, but the mortality rate increases to 50% as they develop septic shock [4]. Mortality risk factors of patients with TB requiring ICU admission have been documented; however, there are no conventions in place to improve mortality of TB-related septic shock. Some of these risk factors that predict in-hospital mortality include acute respiratory failure requiring mechanical ventilation, acute respiratory distress syndrome (ARDS), co-infection with the HIV virus, extensive fibrocavitary disease and consolidation on chest radiographs, and multiple organ failure [5].
It is important to recognize that some cases of TB-related severe sepsis and septic shock will improve with early administration of parenteral anti-TB therapy. Arya et al. described a similar presentation of TB-related septic shock involving a patient with high-mortality risk factors which demonstrated the effectiveness of initiating empiric therapy based on clinical suspicion of TB [6]. It is the high index of suspicion in these cases that likely led to increased patient survival rates. A retrospective cohort study by Hazard et al. concluded that empiric anti-TB therapy was associated with improved survival rates in patients with severe sepsis [7].
A TB infection results from inhalation of airborne particles that are <5μ in diameter, those of which carry few organisms. Upon entering the upper respiratory tract, the droplets deposit into the subpleural airspaces. Many particles are cleared by host alveolar macrophages; however, the TB bacilli can replicate within an activated host macrophage causing a primary infection.
Lipoarabinomannan (LAM), a glycolipid on the mycobacterium cell wall with virulent properties, interacts with carbohydrate recognition domains on the surfaces of macrophages and dendritic cells. By inhibiting SHP-1, a tyrosine-protein phosphatase, mycobacterium may be able to survive in host cells for long periods (Figure 4).
FIGURE 4: Survival of tuberculosis in host cells
Proposed mechanism in which lipoarabinomannan (LAM) directly activates SHP-1 through phosphorylation. SHP-1, a signaling molecule that regulates many different cellular processes, may eventually deactivate proinflammatory cytokines such as IL-12, that plays a role in macrophage activation. This inability to release INF-y and activate macrophages promotes intracellular survival of TB in host cells [8,9].
Without prompt treatment, TB can continue to spread bronchogenically and the caseous material causes the destruction of elastic and muscular components of bronchial walls resulting in widespread bronchiectasis, as seen in our patient.
TB continues to spread both lymphatically and hematogenously causing a systemic response. Proinflammatory mediators promote leukocytosis and stimulate the release of other cytokines, all of which act as pyrogens, activate WBCs and macrophages, promote chemotaxis, cause immunosuppression, and stimulate both coagulation and fibrinolytic activation [10]. Arachidonic acid and adhesion molecules further cause vascular permeability and migration of aforementioned leukocytes. There is a link between tumour necrosis factor-α (TNF-α) and activation of the complement pathway, which further enhances neutrophil trafficking and inflammation [11]. This leads to tissue ischemia, cytopathic injury, and increased programmed cell death, eventually resulting in widespread multi-organ damage [12]. The rearrangement and distribution of intravascular fluid, inhibition of vasopressin release, and upregulation of vasoactive mediators such as nitric oxide results in hypotension or shock, as seen in our patient.
At the level of the lung, endothelial damage in the respiratory vasculature disrupts blood flow and increases microvascular penetrability, leading to interstitial and alveolar pulmonary edema [13]. The edematous fluid protein to plasma protein ratio is nearly 0.95, which is significantly higher than those with cardiogenic pulmonary edema [14]. This proteinaceous fluid destroys pneumocytes which increases surface tension, traps leukocytes, and further damages lung vasculature. Diffuse alveolar damage (DAD) ensues, a significant ventilation/perfusion mismatch occurs, and ARDS develops. When widespread TB-related lung damage occurs, ARDS is hardly reversible, especially when complicated by multi-organ failure.
TB and adrenal insufficiency
TB is one of the most common causes of adrenal insufficiency worldwide and a disseminated TB infection has the potential to spread to the adrenal glands. TB-related adrenal insufficiency may manifest as an acute adrenal crisis, especially in the setting of significant stress caused by a new infection. This is, in part, due to mycobacterial destruction and caseous necrosis of the adrenocortical tissue [15]. We considered our patient's symptoms to be a manifestation of a TB-related adrenal crisis; however, it is uncommon for the adrenal glands to be the only infected organ. Adrenal involvement was found in 6% of patients with active TB in a 28-year-autopsy series, and only one-quarter of these cases listed the adrenals as the only site of infection [16]. Nonetheless, adrenal insufficiency is always important to consider in cases involving disseminated TB. To aid in the diagnosis, one can consider sending a serum cortisol level, adrenocorticotropic hormone (ACTH) stimulation testing, and utilizing an abdominal CT to evaluate for hemorrhage, inflammatory cell infiltration, or adrenal enlargement.
Treatment of TB-related septic shock
Immediate management involves the timely administration of supplemental oxygen and fluids [17]. The administration of balanced crystalloid IV fluids like ringers lactate should be started within the first hour. Two empiric antibiotics intravenous antibiotics from different classes should also be administered to cover for the most common pathogens.
Next, an acid-fast sputum smear and culture should be collected, and laboratory data including complete blood count (CBC), complete metabolic panel (CMP), prothrombin and partial thromboplastin time (PT/PTT), D-dimer, serum lactate, blood cultures, arterial blood gas (ABG), HIV Ag/Ab, betahuman chorionic gonadotropin (B-hCG), and procalcitonin should be obtained. Clinicians should consider starting empiric TB treatment while awaiting acid-fast bacilli (AFB) smear results. If the clinical presentation and imaging findings are consistent with TB and the patient is in septic shock related to TB, anti-TB medications need to be administered as soon as possible. It is important that clinicians do not wait for the sputum results before initiating treatment. With severe sepsis and septic shock, the timing of antibiotics directly affects the outcome. This is also valid for a tuberculous-related septic shock [3]. Further, clinicians must consider a parenteral route for treatment in patients with suspected disseminated TB and septic shock. This is because oral absorption is unpredictable considering poor splanchnic circulation in a state of shock and poor absorption may lead to poor outcomes.
The preferred regimen for the treatment of pulmonary and extrapulmonary manifestations of TB including TB-related sepsis shock without HIV includes two months of RIPE, followed by a continuation phase for an additional four months [18].
Vasopressors are useful in patients who are continuously hypotensive despite adequate fluid management. First-line vasopressors include norepinephrine; second or third agents such as epinephrine, dobutamine, or vasopressin may be used in the setting of reduced cardiac output. Steroids can be considered for six weeks in those with TB meningitis or TB pericarditis, in those with persistent septic shock despite adequate fluid resuscitation, and in those when TB adrenalitis is suspected.
Conclusions
This case highlights the possibility of TB-related septic shock in the critical care setting and the importance of prompt intervention, including treatment with anti-TB agents and therapies to correct hypoxemia and hemodynamic instability. Clinicians should consider administering parenteral therapy as soon as possible while waiting for additional results, which may likely contribute to positive patient outcomes. This case also calls attention to possible mechanisms contributing to TB-related septic shock, as well as acknowledges the link between disseminated TB and adrenal insufficiency.
Additional Information Disclosures
Human subjects: Consent was obtained or waived by all participants in this study. Conflicts of interest: In compliance with the ICMJE uniform disclosure form, all authors declare the following: Payment/services info: All authors have declared that no financial support was received from any organization for the submitted work. Financial relationships: All authors have declared that they have no financial relationships at present or within the previous three years with any organizations that might have an interest in the submitted work. Other relationships: All authors have declared that there are no other relationships or activities that could appear to have influenced the submitted work. | 2,765 | 2022-12-01T00:00:00.000 | [
"Medicine",
"Biology"
] |
Macroeconomic and Idiosyncratic Factors of Non-Performing Loans: Evidence from Pakistan’s Banking Sector
Using panel data approach in the Pakistan banking sector over the period 2010 to 2016, we examine the bank-specific and macroeconomic determinants of non-performing loans. We use quantitative research design with OLS random effect model. Furthermore, we use various regression and correlation analysis in this study. We find that rise in capital adequacy ratio, bank size, GDP growth rate, and inflation, reduces the non-performing loans (NPL) ratio. Our results also show that a rise in loan loss provisions enhances the NPL ratio. Our results suggest that banks with poor asset-quality can sabotage the growth of fiscal as well as the economic sector. Outcomes of the study emphasis on the need to clear-out the NPLs to keep financial sector sound. NPLs can cause high loan loss provisions which affect the capitalization of banks that ultimately impacts fiscal and economic growth. Bank supervisory agencies should therefore pay attention to monitory and macroeconomic policies of the banks. This study examines the impact of idiosyncratic and macroeconomic determinants of non-performing loans on banks’ asset quality using recent data from 2010 to 2016, the time period when major banking sector reforms were launched.
Introduction
Among all the financial institutions, the role of banks is most significant and distributional. The bank is a body which amasses deposits from regulars and gives loans to organizations and indi-restricted with the rate of SBP (State Bank of Pakistan, 2017). Since the last two decades, the banking system in Pakistan has been well established because of a series of liberalization in policies and financial reforms. The progress of the banking industry is due to the vigilant supervision of the State Bank of Pakistan. 82% of the financial sector of Pakistan comprises banks that are categorized further as conventional, Islamic, specialized and foreign banks (State Bank of Pakistan, 2016).
The banking sector, not only in emerging but also in mature economies, observes many problems. The poor performance of banks results from many factors such as lack of management efficiency, low capital adequacy ratio, and poor asset quality. One of the biggest problems of the banking sector is the non-performing asset (Sharma, Tiwari & Sood, 2013). Commercial banks try to invest as much as possible in the form of loans and credit for the maximization of profit which shows most of the assets of banks exist in the form of loans but there is a huge risk of debt recovery (Achou & Tenguh, 2008). Although the loans are the largest assets of the banks and a major source of income, there is great risk in granting loans (Casu & Girardone, 2006;Honey, Tashfeen, Farid & Sadiq, 2019).
A huge amount of non-performing loans can influence the intermediary role of banks for the progress of the economy and nation. Research practices show that non-performing loans are top indicators of financial crises (Brownbridge, 1998;Greenidge & Grosvenor, 2010), however, poor and inefficient management and inefficiency of firms are also vital factors for non-performing loans (Fan & Shaffer, 2004;Girardone, Molyneux & Gardener, 2004). Failure to repay the debts causes the emergence of non-performing loans, which is the greatest financial problem (Heffernan, 2005). According to IMF (2009) definition: " A loan is non-performing when payment of interest and principal are past due by 90 days or more, or at least 90 days of interest payments has been capitalized, refinanced or delayed by an agreement or payments are less than 90 days overdue, but there are other good reasons to doubt that payment will be made in full." Empirically, the occurrence of banking crises is closely related to a huge accumulation of non-performing loans that contains a major share of assets of an insolvent bank. Association of nonperforming loans and banking crises can be proved from different financial crises in the world such as Asian financial crises of 1997 which spoiled the financial system and economies of many countries, in Indonesia 60 banks were collapsed and their 75% loan portfolio became non-performing, financial crises of 2007-2008 in America which then ruled over in different countries and cause financial instability (Caprio & Klingebiel, 2002). Non-performing loans agitate the overall bank efficiency and the high level of nonperforming loans depicts the huge amount of credit defaults. The growth of non-performing loans involves the necessity of provision, which eventually decreases the profit level. The branch manager should know the causes of bad loans and should verify the customers before providing the loans because an effective and efficient monitoring system can increase the performance of the banking system which ultimately has a positive impact on the economic growth (Sharma et al., 2013).
Nigeria's banking industry observed a sharp upswing in the ratio of non-performing loans (NPLs) by 220% from December 2015 to December 2016 as the amount of NPLs climbed from 0.65 trillion to 2.08 trillion. NPL to total loan ratio (NPL ratio) increased from 4.88% to 12.80% in one year resulting in decreased profitability of commercial banks by 30.16% (NDIC). Nonperforming loans were the major cause in Nigeria which limits the segmental growth of the economy (Boudriga, Taktak & Jellouli, 2010;Adeyemi, 2011;Bebeji, 2013).
Problem Identification
Non-performing loans are closely related to banking crises (Kroszner, Laeven & Klingebiel, 2007) as non-performing loans are important indicators of financial stability and in an increase in the level of non-performing loans cause bank failure (Bardhan & Mukherjee, 2016;Ghosh, 2015;Kasman & Kasman, 2015;Nkusu, 2011).In 2006, the level of non-performing loans started to increase in America which lead to the subprime mortgage crash in 2007 (Greenidge & Grosvenor, 2010).Global financial crisis of Volume 1 Issue 2; August 2019 2007-2009 which damaged the USA economy and economies of many countries was also because of the non-performing loan (Adebola, Yusoff & Dahalan, 2011).
An unparalleled climb of non-performing loans in the Japanese banking sector during the 1990s generated a protracted economic collapse, during the chaos government undertook the stabilization arrangements by advancing insurance, inserting public capital and bailing out concerned banks which results in a decrease of government assets (Hoshi & Kashyap, 2010;Montgomery & Shimizutani, 2009). Credit crises in Mexico after 1995 were also due to the bad loans because financial institutions were loaded with a huge amount of credit with the negative value which decreases their capability to provide further loans to different sectors of the economy (Krueger & Tornell, 1999).
A study on commercial banks of Bangladesh shows that managing non-performing loans are important in developing investor confidence. If their volume is not monitored appropriately, it may harm the opportunities for new borrowers. The volume of default loans of banks listed on Dhaka stock exchange has been increasing at a shocking rate and this situation is due to excessive political and illegal interference. The amount of non-performing loans was Tk.546.57 billion till 2015 which was Tk. 427.3 billion in 2012 and Tk. 200.1 billion in 2006, so high volume of nonperforming loans cannot be profitable for an economy because non-recovery of funds confines the re-use of funds which leads to an economic sluggishness (Haruna, 2013;Buchory, 2015). In Pakistan, 80% of the banking sector is privately owned and when private banks are not willing to disburse loans to investors, it results in interest rate increase and diminishes the profitability of the banking zone that exemplifies the weak state of the economy (State Bank of Pakistan, 2016). As shown in Figure 1, Pakistan has the 3rd highest percent of non-performing loans in South Asia which is an alarming situation for the economy. An increase in non-performing loans would logically decrease the worth of assets, which subsequently leads to extensive losses and significant retrenchment in obligatory capital. The swift climbs in non-performing loans borders the lending activities of banks which ultimately has consequences in rendering economic proceedings due to stumpy speculation of money and reflected as a sign of financial crises (State bank of Pakistan Working Papers, 2015). In Pakistan, NPLs are also affecting the economic and financial sector performance. Despite the efforts of the Central Bank to control the ratio of NPLs the performance figures of the last 25 years didn't fall double-digit (SBP, 2016). In Pakistan from 1995 to 2016 the average level of NPLs is14.87% which is alarming for financial sector growth.Pakistan is 24 th in terms of the highest level of NPLs states (State bank of Pakistan 2016).
Problem Statement
Boudriga et al., (2009) general perspective and also for lending operations but still NPLs are a big problem for local and international regulators. The negative impact of non-performing loans on banks and the economy is a problematic issue for supervisory institutions and policymakers all over the world (Sočuvková, 2013).
Literature Review
The financial stability of the economy and its growth are considerably influenced by the level of non-performing loans. An increased level of the non-performing loan is the symbol of shrinkage of economic progress due to the non-performance of assets that cause a high rate of unemployment and a gradual decrease in asset prices. (Klein, 2013;Farhan, Sattar, Chaudhry & Khalil, 2012;Nkusu, 2011;Sapkota, 2012). The association between the non-performing loans and idiosyncratic and macroeconomic factors is extensively investigated in the current literature due to the significant impact that non-performing loans have not only on financial institutions but the economy as well (Reinhart & Rogoff, 2011;Castro, 2013;Makri, Tsagkanos & Bellas, 2014;Chaibi & Ftiti, 2015;Saba, Kouser & Azeem, 2012). Changes in the macroeconomic condition of a country lead to a change in the lending practices and their utilization as unemployment and rate of interest have a significant impact on loan quality of banks. Many prevailing studies explore macroeconomic determinants of NPLs for different countries, most of the studies find out the invers connection among macroeconomic atmosphere and nonperforming loans. Rate of inflation, unemployment, external debt to GDP growth rate, amount of loan, credit to the private sector, exchange rate, share price, and lending rate of interest are the indicators of the non-performing loans and have a substantial impact on the economic growth of the country (Ghosh, 2015;Škarica, 2014;Zeng, 2012;Espinoza & Prasad, 2010;Dash & Kabra, 2010;Swamy, 2012).
Studies on the banking sector which explore the impact of macroeconomic aspects on the level of non-performing loans show that GDP growth rate, rate of inflation, rate of interest, and exchange rate has a negative effect on non-performing loans in longterm perspective while lending rate of interest is positively related with non-performing loans, as an increase in lending rate leads to decrease in reimbursement ability of borrower because it also increases the rate of inflation which reduces the monetary value of currency and has a negative impact on non-performing loans (Badar & Javid, 2013; Warue, 2013). Chiorazzo, D'Apice, Morelli & Puopolo, (2017) conclude that GDP growth rate, a high rate of interest, and efficient judicial system are major macroeconomic determinants of non-performing loans which influence the payback capacity of the borrower.
Empirical studies show that bank-specific factors such as last year NPLs ratio, bank size, net interest margin, credit risk, liquidity, ownership structure, corporate governance, legal terms of the loan agreement, and the current rate of loan growth have significant impact on the volume of non-performing loans. Macroeconomic factors such as inflation in previous as well as the current year, GDP per capita growth and exchange rate, interest rate, and inflation enhance the non-performing loan volume. However, in large banks, both types of factors, bank-specific and macroeconomic, influence the non-performing loan ratio while in small banks non-performing loans have only influenced by bank-specific factors. (Amuakwa & Boakye, 2015;Klein, 2013;Inekwe, 2013;Dash & Kabra, 2010;Swamy, 2012;Sadiq, et al., 2017).
A study by Farhan et al., (2012) on Pakistani banking sector shows that interest rate, energy crises, inflation, unemployment, and exchange rate have a significant positive impact on nonperforming loans of banks while GDP growth has a negative impact on non-performing loans ratio, this study also shows how term loans become bad loans due to low production of industrial sector because of energy crises. Anisa (2015) state that deposit rate, loan to deposit ratio, and the lending interest rate have a positive impact on non-performing loans while solvency ratio of bank and GDP growth rate have a negative impact on non-performing loans. Angelos, Louzis, Vouldis and Metaxas (2012) evaluate the Greece banking system and conclude that macro-economic factors such as GDP, exchange rate, unemployment, and bank-related factors possess the ability to influence the level of non-performing loans of each category such as corporate loans, house loan, and car loans, etc.
The diverse trend towards the association between GDP growth rate and the magnitude of non-performing loans has been observed in the literature. GDP and non-performing loans are positively interlinked in a few studies, though frequent studies also show the negative correlation between non-performing loans and GDP. GDP growth rate for the same period has a negative effect on non-performing loans while the latency GDP growth rate has a positive effect on non-performing loans. Since GDP increases its indicators to a higher level of income which boosts the capability of borrowers to reimburse loans. When there is a depression in the economy (slowed or negative growth of GDP) the level of bad obligations will upturn (Salas & Saurina, 2002;Khemraj & Pasha, 2009;Dash & Kabra, 2010;Shingjergji, 2013).
Macroeconomic factors have an immense impact on the profitability of banks because these factors are not in the control of banks and management due to their impact at the macro level, so they influence the different levels of growth according to the size and nature of bank. Deterioration in the economic condition of a country reduces the debtor's ability for repayments because it decreases the per capita income (Mileris, 2014). Inflation is also assessed as the significant macroeconomic determinant of nonperforming loans, although its relation is inconclusive. Loan payment capacity can be affected by the inflation positively as well as negatively depending upon the situation of the economy, as a high rate of inflation will lead to a decrease in the capacity of the borrower for loans because the monetary value of his income will decrease by the decrease in the value of the currency. Inflation rate has a positive relation with non-performing loan as a lower rate of inflation has a significant positive impact on financial condition of the borrower and thus on its repayment capacity (Mileris, 2012;Khemraj & Pasha, 2009;Gunsel, 2012;Thiagarajan & Ramachandran, 2011;Abid, Ouertani & Zouari-Ghorbel, 2014) while inflation has a negative association with non-performing loans according to (Warue, 2013;Shingjergji, 2013).
When the rate of interest is high then, the organizations generate a high rate of return to cover the cost of capital to avoid the insolvency element. The higher the interest rate increases the debt burden which declines repayment capacity of the borrower and ultimately the size of the non-performing loans increases (Aver, 2008;Castro, 2013;Skarica, 2014;Ghosh, 2015;Curak et al., 2012;Bardan & Mukjerjee, 2016).
Rate of unemployment has a positive relationship with the NPLs as an increase in unemployment leads to a decrease of a debtor's income, which disturbs their capability to reimburse the loan. Deviations in unemployment are reflected as a good sign of the recession (Charalambakis, Dendramis & Tzavalis, 2017). Increase in the rate of interest leads to a higher rate of the unemployment which has an ultimate impact on non-performing loans because unemployment reduces the flow of cash of a household which decreases the consumption in economy; on the other side, an increase in unemployment rate also affects the firm's cash flow, and it results in a decrease in their production. Makri et al. 2014;Chaibi, Hasna & Fititi, 2015). Furthermore, non-performing loans are also positively related with expected lending interest rates while the rate of interest has a negative link with the level of non-performing loans, because increase in rate of interest climbs the rate of inflation which decreases the purchasing power and thus the repayment capacity of borrower get decreased due to unemployment (Ali, Shingjerji & Iva, 2013;Akinlo & Emmanuel, 2014;Vardar, Gulin & Ozguler, 2015;Messai & Jouini, 2013;Skarica & Bruna 2014;Donath et al., 2014). Ahmad and Ariff (2007) state that the credit risk is the most harmful one among all the risks the bank face, as non-performing loans affect the bank profitability and long-term operations. The high volume of problem loans in the credit folder of banks is incompatible to banks in attaining their goals. Adebola et al., (2011) state that the high build of non-performing loans indicates financial instability of the bank. Garr (2013) discusse that the credit risk strategy of a bank is contingent on the economic condition and its management is multifarious due to the fickle nature of macroeconomic dynamics and bank-specific features. Credit risk management is an important factor to determine the financial performance of banks because effective credit risk management leads to the greater financial performance of the banks and their profitability (Alshatti & Sulieman, 2015;Gizaw, Kebede & Selvaraj, 2015). An increase in the default rate damages the entire banking system and as a result, the inflation, rate of interest, stock index, and industrial outcomes are affected by these defaults (Boss, 2002).
Bank size is a significant factor for non-performing loans. There are mixed results of the studies on the consequence of bank size on the level of non-performing. An inverse connection is endorsed to the point that large banks have better risk supervision tactics to come up with issues of non-performing loans. (Rajan, Rajiv & Dhal, 2003;Sales & Saurina, 2002). Large banks have a better opportunities to deal with non-performing loans, so they have a low level of non-performing loans hence, they find a negative relationship between bank size and non-performing loans (Hu Li & Chiu, 2004;Swamy, 2012) conduct a research study in Nigeria for 20 years and the results show that huge ratio of non-performing loans reduce the performance of banks as it reduce, the return on capital employed in both short run and long run. H4: Capital adequacy has a significant impact on non-performing loans. H5: Credit Risk has a significant impact on non-performing loans.
Methodology
The current study opts to analyze the determinants of nonperforming loans and their impact. The research problem is premeditated by the use of descriptive and explanatory research design which is concerned with outcomes of what and how the phenomenon has occurred. The descriptive research design allows for greater generalizability of the findings (Gremi, 2013;Park & Zhang, 2012;Mileris, 2012;Castro & Vitor, 2013;Igan, Deniz & Pinheiro, 2011;Vogiazas, Sofoklis D & Nikolaidou, 2011;Salas & Saurina, 2002). The explanatory research design describes the cause-and-effect relationship between the dependent and independent variable which is also observed in the current study in various macroeconomic variables (Kothari & Rajagopalachari, 2004).
Sampling
The way through which we select our sample is called a sampling technique. The target population for this study is all the commercial banks of Pakistan which include foreign banks, private banks, and public banks, which are registered with the State bank of Pakistan. Data sampling is done by using stratified sampling, stratas are made according to the ratings of the banks issued by PACRA. Banks that have ratings of AAA, AA+, AA-and AA are selected. Fourteen banks fall in this category of Rating, so these 14 banks are sample of the current study. Panel data of selected commercial banks in Pakistan covering the period from 2010 to 2016 is studied. The use of panel data instead of cross-sectional or time series is very beneficial in terms of efficiency of econometrics estimates because it contains a large number of observations which leads to the higher number of degree of freedom that helps in finding the answer of the wide range of questions (Hsiao & Cheng, 2014).
Data Collection
The secondary data for the study is collected from, Annual statements of banks World Bank annual database created by the world bank
Model Estimation
To analyze the bank-specific and macroeconomic determinants of non-performing loans, the following equation is supposed: NPLi,t= βo + β1GDPit + β2INFit + β3BSit + β4CARit +β5CRit +Ɛi,t NPLi,t = NPL ratio of bank i at time T β1GDPit = GDP growth rate at time T β2INFit = INF rate at time T β3BSit = BS at time T β4CARit = CAR at time T β5CRit = CR at time T Ɛi,t = error term Non-performing loans ratio has a minimum value of 1.40 and the maximum value of 32.8 with the mean value of 12.09 showing the deviation of 5.64% from its mean value. This shows that selected sample banks incurred 12.09% non-performing loans on average from its total loans. Credit risk which is measured as loan loss provision ratio in this study has a range from 0.02 to 6.02 with a mean value of 0.62 shows the deviation of -0.192 from its mean value.
Descriptive Statistics
Capital Adequacy ratio has a minimum value of 1.05 and a maximum of 49.7 % with a mean value of 15.7% and has a deviation of 9.28%. The mean value for selected sample banks shows that CAR is higher than the minimum requirement of CAR according to the state bank of Pakistan which is 10%.
Bank size has a value of range from 5.03 to 6.5 with a mean value of 5.68 which shows the highest standard deviation of 5.35 that shows the presence of high variation in terms of size in selected banks. GDP growth rate has a range from 2.58% to 5.74% shows the mean value of 4.17 and has a standard deviation of 3.25. Inflation has value in the range from 2.5% to 13.9% and a mean value of 8.04% shows that it is deviated from mean value by 3.82%.
Gross Domestic Product
GDP has a negative/inverse relation with non-performing loans in the banking sector of Pakistan according to the results. A coefficient estimate of Gross domestic product is -1.24 which shows a negative relation at a 95% confidence interval. This inverse relation depicts that an increase in growth rate would decrease in the amount of non-performing loans which is based on the fact that a good growth rate shows the good health of the economy and ultimately the standard of living of the people. An increase in growth rate shows that the economy will perform its best level and hence the standard of living is going to be enhanced. GDP growth rate shows an increase of 122% from 2010 to 2016.This growth also shows that people of the country held good economic status both in terms of individual living and business entities, which prevents people from being a defaulter of their loans as NPLs ratio decreases to 31% from 2010 to 2016. Consistent results have been found in the studies of (Fofack & Hippolyte, 2005;Saba et al., 2012;Louzis et al., 2010;Klien & Nir, 2013). Graph 1 shows that in 2010 the GDP growth rate is 2.58% and the percentage of NPLs is 14.75% which tends to decrease at the rate of 10.06% with the increase in GDP growth rate.
Inflation
Another macroeconomic variable which is considered in this study is inflation. Inflation has a negative relation with non-performing loans in this study which is based on the notion that inflation reduces the time value of money because too much money chases too few goods. As it also affects the value of the remaining debts so the borrower will feel easy to repay his debts. Results show that inflation has a coefficient estimate of -0.09 and has a significant relation at a 90% confidence interval. Results of this study are supported by (Khemraj, Tarron & Pasha, 2009;Warue, 2013: Shingjergji, 2013) while few studies have also shown a positive relationship with non-performing loans such as (Nkusu, 2011;Farhan et al., 2012). Inverse relation of inflation and non-performing loans can be seen in Graph 2 as in 2015 the percentage of NPLs is 11.36% which declined to 10.06% in 2016 with the increase in inflation from 2.5 % to 3.8%. Graph 2: Inflation Rate vs NPLs
Capital Adequacy Ratio
Capital adequacy has a minimum requirement of 10% according to the Prudential Regulation of the State Bank of Pakistan. CAR is a bank-specific variable in this study which has a negative coefficient of -0.344 and it is statistically significant at a 99% confidence level. Results show that CAR has an inverse relation with nonperforming loans which is supported by the justification that wellcapitalized banks would sustain the different types of risk and losses arises from them because this enough capital would lead its better regulation process. As minimum CAR of 10 % by Prudential regulation of State Bank of Pakistan is maintained by most banks in the banking sector of Pakistan so they exhibit negative trends towards non-performing loans. Similar results are found in the studies of Zhang and Shihong (2012), Swamy (2012) and Makri Tsagkanos & Bellas (2014).
Bank Size
Researches have shown both positive and negative relation of the size of a firm with non-performing loans. Few studies exhibit that bank size has a direct and positive relation towards bad loans means larger the size of banks higher will be the ratio of their nonperforming loans which is clinched by the fact that larger banks may avoid over-monitoring of borrowers not only after advancing the loans but also before advancing the loans. The problem of dis- torted information such as lack of disclosing about financial status in larger banks would also be fundamental to a rise in the level of problematic loans. Notwithstanding the above phenomena, bank size in this study has a negative coefficient and significant at a 99% confidence level. This relation illustrates that increase in the size of a bank would decrease their volume of non-performing loans in Pakistani banking sector because bigger size banks have better monitoring system not only after advancing the loans that where are these loans being used and what is the purpose of taking loans but larger banks also have a monitoring system of the background of the loan taking firms and individuals. Bigger banks have an efficient and effective risk management system and better information system that how they would maintain equilibrium to minimize the risk of defaults (Al-Smadi, Mohammad & Ahmad, 2009;Godlewski, 2005).
Credit risk
Credit risk, which is measured in this study as loan loss provision, has a significant positive relationship with non-performing loans in this study. The positive result illustrates that high loan loss provision depicts that banks face a high level of non-performing loans. This result shows that banks owe a high amount of provision because of the perceptions that customers will not able to pay off their loans. Moreover, poor credit quality is also an issue that increases the risk portfolio of banks. P-value shows that this positive relation is confirmed on a 95% confidence interval. Results of this study are aligned with the results of (Chaibi & Ftiti, 2015;Boudriga, Boulila & Jellouli, 2009;Messai & Jouini, 2013)
Conclusion and Recommendation
The major objectives of this study are to examine the impact of different macroeconomic determinants of non-performing loans in the banking sector of Pakistan. To attain this goal, the quantitative research approach is used along with panel data analysis on the period from 2010 to 2016. Random effect model has been used for the analysis of the data. To accomplish the analysis, EVIEWS version 10 is used. GDP rate is negatively and statistically significant which shows that whenever the economy will be on its peak, the value of cash held for household and business will increase which will reduce the behavior of nonpayment of their financial obligations. Inflation is also a significant element of non-performing loans and has a negative impact on bad loans because the increase in the rate of inflation decreases the worth of cash, therefore, it becomes easy to oblige financial obligations as the value of outstanding loans to become less. So overcoming the impact of inflation, there should always be a moderate level of inflation not very low or high. Bank size and credit risk also have a significant impact on bad loans. Banks should pay attention to their lending policies and monitoring system to avoid problem loans. Moreover, wellcapitalized banks do not face the problems of bad loans.
Recommendation
According to the findings of this study, GDP growth rate, inflation, bank size, capital adequacy ratio, and credit risk have a significant impact on non-performing loans. To evaluate the macroeconomic impact, the concerned authority should make effective macroeconomic policies to avoid the problem of bad loans while to cure the problem of bad loans, better risk management systems, better lending policies and efficient monitoring systems of the borrower with a check of symmetry of information should be followed. Vigilant and vibrant credit policies would incorporate appropriate customer selection and sanction processes with clear retrieval policies. To ensure a sound financial system, the State Bank of Pakistan should direct the commercial banks that credit facility to a potential borrower would not be granted without the prior written approval of the State bank of Pakistan. Moreover, commercial banks should pay their attention to modern and inventive means of increasing their interior financial capability so they can handle their financial matters efficiently.
Limitations of the Study
This study considered only fourteen banks for a seven-year periods with three banks specific and two macroeconomic variables.
Direction for Future Research
As this study has considered only two macroeconomic factors of non-performing loans in the banking sector of Pakistan, and in this econometric model all macroeconomic determinants are not in- | 6,830.6 | 2019-08-30T00:00:00.000 | [
"Economics"
] |
Conducting polymers as electron glasses: surface charge domains and slow relaxation
The surface potential of conducting polymers has been studied with scanning Kelvin probe microscopy. The results show that this technique can become an excellent tool to really ‘see’ interesting surface charge interaction effects at the nanoscale. The electron glass model, which assumes that charges are localized by the disorder and that interactions between them are relevant, is employed to understand the complex behavior of conducting polymers. At equilibrium, we find surface potential domains with a typical lateral size of 50 nm, basically uncorrelated with the topography and strongly fluctuating in time. These fluctuations are about three times larger than thermal energy. The charge dynamics is characterized by an exponentially broad time distribution. When the conducting polymers are excited with light the surface potential relaxes logarithmically with time, as usually observed in electron glasses. In addition, the relaxation for different illumination times can be scaled within the full aging model.
Scientific RepoRts | 6:21647 | DOI: 10.1038/srep21647 between carriers and a Coulomb gap in the single-particle density of states 19,20 . If the radiation frequency exceeds the Coulomb gap energy, charges are randomized and reestablishing the equilibrium configuration is usually a very slow process.
Slow logarithmic relaxation of the electrochemical doping potential was observed over nine orders of magnitude in polymeric materials 21 . A dependence of the logarithmic shift with scan rate and aging phenomena were also observed. The universal features of slow relaxation were attributed to a hierarchical series of processes, following ideas from kinetic studies of the decay of persistent photoconductivity in semiconductor structures 22,23 .
Slow relaxation in the photoconductivity in poly(phenylenevinylene) (PPV) films was studied by Lee et al. 24 . A roughly logarithmic behavior was observed, as well as an ω 0.66 dependence of the photoconductivity with chopping frequency, which was explained again in terms of hierarchical processes working in series. Logarithmic relaxation over many orders of magnitude was observed in photoinduced conductivity in organic field effect transistor 25 . Under illumination, the photoinduced excess current also increases logarithmically.
All experiments mentioned describe conductivity at the macroscopic scale. Some mesoscopic properties have also been analyzed 13,26,27 , but the nanoscale has never been reached in electron glass studies. The scanning force microscopy (SFM) technique offers the double opportunity of being able to observe nanoscopic properties of interacting systems and to follow the relaxation behavior for a quantity other than conductivity. In this letter, we use the scanning Kelvin probe microscopy (SKPM) to measure the surface potential (SP) in order to explore a possible, very natural explanation for slow relaxation in conducting polymers based on electron glass ideas. Firstly, we study the SP distribution in equilibrium and see a well defined domain structure with spatial and temporal correlations fully compatible with the electron glass model. Secondly, the sample is excited by irradiation with green light for a certain time and we monitor how the sample SP relaxes to equilibrium once the excitation is switched off. We have observed a logarithmic behavior characteristic of glassy interacting systems.
Results
Nanoscale charge domains. The nanoscale SP of MEH-PPV thin films has been studied by means of SFM movies (see Materials and Methods). These movies are a powerful tool to investigate dynamic properties at high resolution and have been previously used to study, for example, single atom diffusion 28 or phase transitions 29 .
Here we utilized them to study the dynamics of electrons on the surface of thin polymer films.
Domain characterization. Figure 1A,B show a typical topographic and SP frame, respectively, extracted from a movie (20 min/frame)(Video S1 and S2 Supplentary Information). The features observed in the topographic and the SP images are essentially independent, as can be seen directly or deduced from the lack of a central peak in their cross-correlation image (inset between the two figures). The roughness of the topographic image is about 1 nm and the typical lateral size of the topographic features is about 40 nm. The histogram of the SP (black curve in Fig. 1E) shows a Gaussian distribution, with a standard deviation of σ S = 71 mV. The mean SP domain size is about 50 nm. It has been checked that all frames have the same statistical properties. In particular, σ S is about 70 ± 2 mV in all frames, and surprisingly about three times larger than the thermal energy kT = 25 meV at room temperature.
To further explore the properties of the films, the average of all (forward) topography and SP (S av (x, y)) frames of a movie (Fig. 1C,D respectively)has been calculated using Eq. (3). The morphology of the sample barely changes during the whole movie. On the contrary, the SP domains change appreciably from frame to frame. This change is not due to instrumental noise (≈ 10 mV), as proved by the fact that the forward and backward scan directions basically coincide ( Fig. S1 Suplementary Information). In addition, the average SP image presents a stable and well-defined structure, implying that it is not due to random fluctuations, but to real electronic properties of the sample. The histogram of the average SP, S av (x, y), is shown in Fig. 1E (red curve). The standard deviation of S av (x, y) is 14 mV, and thus much smaller than the standard deviation of single frames (70 mV). Therefore, the average charge of a domain must be much smaller than its typical instantaneous charge. The cross correlation between the average topography and the average SP (inset between Fig. 1C,D) shows a small peak in the center (of height 1/20) indicating a weak correlation between the average SP and the topography, which can hardly be resolved for individual frames.
The standard deviation over time δ S (x, y) (Eq. 4) is shown in Fig. 1G. δ S (x, y) is essentially constant (≈ 70 mV) and shows no significant spatial structures. This constant is quite large, showing that the SP has a high variability over time independent of x and y. Interestingly, this value is similar to σ S in a typical single frame.
Time correlations. In order to study dynamic properties we compare successive movie frames through their cross-correlation. Frames in Fig. 1 were taken at a rate of 20 minutes/frame to minimize noise and the cross-correlation between consecutive frames is basically zero. To increase time resolution, fast images were acquired (from 15 to 120 seconds/frame). Figure 2A-D shows four representative frames extracted from a fast movie (60 s/frame) corresponding to different times (Video S4 Suplentary information). The autocorrelation of the first frame (t = 0) together with cross-correlation between the first frame and successive frames ( Fig. 2A-D bottom panel) show that the cross-correlation decays over a few frames, corresponding to a characteristic decay time of about three minutes.
To analyze further the dynamical response of the system, we measure the SP evolution with time at a fixed point. In this way we are able to explore time scales from tens of milliseconds (limited by the response of the microscope) to several minutes (limited by the drift of the apparatus, which cannot be corrected for a fixed point). From these data, we calculate the time correlation function C(t), (Eq. 4), plotted in Fig. 2E. One can appreciate that the scale of characteristic times is very wide, spanning more than four decades. This wide spread of characteristic times is usually found in glasses and is ascribed to the existence of a hierarchical set of processes. To confirm that the overall behavior at different timescales is consistent, C(t) obtained for two movies acquired at different speeds (Video S4 and S6) has been included ( Fig. 2E green and blue symbols). The agreement is good taking into account that the data were acquired at different speeds and on different samples. The fast decay at short time scales is related to the dynamical response of the apparatus.
Spatial correlations. The study of spatial correlations is a complicated problem since dynamic time scales are comparable or even faster than data acquisition times. Very fast scans are crucial in order to minimize the effects due to temporal evolutions. To overcome this difficulty as much as possible, the SP was measured along the same line for a series of scans. An example is shown in the left inset of Fig. 3, where the horizontal axis is the position along the line and the vertical axis time. One can appreciate vertical bright and dark lines, corresponding to charge domains that are stable on the timescale of seconds to minutes. It can be clearly seen that most domains change their state and some of them change rather suddenly at a given instant probably due to charge jumps. The autocorrelation function for each horizontal line of the SP data is shown in the right inset of Fig Physical interpretation in terms of an electron glass. The first point that should be addressed is the size of the SP spatial fluctuations (σ S = 71 mV), roughly a factor 3 larger than the thermal energy. This can be naturally explained within the electron glass model, since in this model site energy fluctuations are much larger than kT 30,31 . Transition energies are of the order of kT, and are equal to the difference in site energy minus the interaction energy, which can be quite large. Site energy fluctuations can be appreciably larger than transition energies, thus larger than kT. The variation of the SP may be due to disorder energy of the intrinsically disordered polymeric material, to possible defects/trap sites as well as to a different electrochemical potential of crystalline/amorphous regions of the semiconducting polymeric material and, finally, to the Coulomb energy E coul of a charge confined to a region of radius R. Assuming that domains are singly charged, for a radius R = 25 nm the Coulomb energy is 60 meV, quite compatible with the observed variations of the SP.
The spatial fluctuations of the time-averaged SP image, around 14 meV, are much smaller than single frame fluctuations, which indicates that the fluctuations of the average SP are mainly due to changes in the local contact/ chemical potential. Each frame corresponds to a particular arrangement of charges. These charges evolve in a disordered mean effective energy landscape generated by local variation of the material properties and/or some mean effective electron configuration.
What is quite remarkable -in our opinion -is that some electronic events occur on very large timescales of up to minutes, and are thus observable by SFM techniques. The electron glass model is an excellent candidate to explain this, since frustration and interactions can produce characteristic times exponentially distributed. The measurement time of SFM is of the order of a millisecond, therefore any dynamics faster than this time is averaged. In this sense, the observed domain radius should be considered as an upper bound of the real localization length of the charges. Nevertheless, this does not invalidate the previous argument relating domain size and SP fluctuations, since both quantities will be averaged in a similar way.
In order to get further insight into the charge domain dynamics, we have performed the same experiment in a MEH-PPV sample that was previously irradiated with blue light for a short time period (see Supporting Information for details). It is well known that blue light photo-induces MEH-PPV degradation, decreasing the π-bond conjugation length and therefore the material conductance 32,33 . Comparing this degraded sample with the non-degraded one, it is found that although in the degraded sample a charge domain dynamics is still present, the corresponding changes occur less frequently. The decay time of the cross-correlation C(t) increases from about 3 min to about 15 min (see Fig. 2 of Supporting Information). Those results fully support the idea that we are observing hopping dynamics, and that a decrease in conductance is correlated with slower domain dynamics.
The existence of a region where the autocorrelation function is negative is another indication of the importance of interactions, since a region with a given charge is more likely to be near a region with opposite charge. This supports the applicability of the electron glass model to MEH-PPV, which is strongly reinforced by the relaxation experiments that we report below.
Finally, we recall that the SFM technique probes mainly surface properties, therefore the images discussed reflect electron dynamics within a relatively thin surface layer. From our experiments we cannot conclude if the phenomena that we are observing is a pure 2D surface effect or a surface projection of bulk effects. In any case, the dimensionality of the electron glass does not affect the interpretation and relevance of our results.
Slow relaxation.
A suitable tool to study complex systems is to drive them out of equilibrium and to study how they return to equilibrium. To achieve a proper sample excitation the "two pass" method 34 is used. In this protocol the sample is illuminated with several "on"/"off " cycles until the light is completely switched off and the system is allowed to relax. By recording the SP, S(t), along the whole experiment, relevant information from the excited state as well as from the relaxation processes can be inferred. Figure 4A shows the time evolution of S av (t) in a typical MEH-PPV thin film excitation experiment. The horizontal black line shows the average SP of a sample region that has not been previously illuminated. This average SP value characterizes the initial state and will be used as reference to evaluate photo-induced excitations. As shown in the inset of Fig. 4A, which is an enlargement of the main figure, during the "two pass" excitation protocol, the light is switched on and off alternatively in periods of 1.5 s. Green and black dots represent SP values obtained under illumination and in darkness, respectively. It should be noticed that processes faster than the scanning time of a line cannot be resolved. As shown in ref. 34, it is observed that the average S(t) increases when the light is "on", as expected for a p-type material 35 . However, S(t), instead of returning to its equilibrium value during the "off " part of the illumination protocol, decreases to smaller values. As more light excitation cycles are applied, the SP for both with and without illumination periods decreases slowly. This indicates that relaxation in the "off " part of each cycle is not complete and the system experiences a higher and higher degree of excitation as more cycles are applied.
Once the light is definitively switched off, the SP slowly tends (relaxes) to its initial value. This relaxation behavior is qualitatively similar (with opposite sign) to the excitation curve, a feature usually found in electronic glasses 10,14,36 . To further analyze the relaxation tendency, data has been plotted in Fig. 4B on a logarithmic scale with the origin of time taken at the instant at which light is completly switched off. Relaxation of the SP to its equilibrium value is roughly logarithmic over four decades of time. Hence, we are clearly dealing with a very slow relaxation process. The rate of relaxation decreases when the intensity of the light increases. A similar effect occurs in GeSbTe in the presence of persistent photoconductivity 17 , and is associated to an increase in carrier density.
The logarithmic relaxation of the average SP implies the existence of an exponentially broad distribution of relaxation times 37,38 . This type of distributions is obtained whenever the transition rates are exponential functions of smoothly distributed random variables, as naturally occurs in hopping systems 39 . In these systems, states are localized and the transition rates for one-electron hops are given by 9 where τ 0 is a typical phonon time, r the hopping distance, ξ the localization length, and Δ E the transition energy. For interacting systems, simultaneous many-electron hops are possible and this equation is still a valid approximation considering r as the sum of all the hopping distances. Previously, slow relaxation in organic polymers was explained in terms of hierarchical processes with an increasing spatial separation between the negative charges at the surface and the positive charges in the bulk 25 . This model cannot explain the rich variety of phenomena that we see in our systems. We think that the existence of the domain structure of the SP together with the logarithmic relaxation favor an explanation in terms of electron glasses. Slow logarithmic relaxation phenomena have been observed in a great variety of systems, grouped under the name electron glasses, as they all conduct by hopping and interactions between carriers are believed to be important. The details of the mechanism responsible for slow relaxation in electron glasses are not known. Nevertheless, there is a growing consensus that energy relaxation requires a hierarchical series of processes involving an increasing number of particles participating in a hop 40 , whose characteristic times grow exponentially with this number. The existence of SP domains ensure the presence of charges and the fact that they fluctuate in time proves the relevance of hopping processes. Additionally we argue that average SP should be correlated with the macroscopic conductance, reinforcing our idea that we are studying the same phenomena that has been observed in electron glasses. After illumination, in a highly excited state, there are many holes below the chemical potential as well as many electrons above it. This implies that there is a shift in the chemical potential with respect to equilibrium (as long as the single-particle density of states is asymmetric) and, at the same time, there is a increase in conductance because there are many more possible excitations. As the system relaxes, changes in the conductance should be correlated to changes in the average SP. Our observations and the presence of aging (to be discussed in the next section) are difficult to explain with previous models of slow relaxation in organic polymers, while they are naturally explained with the electron glass model.
Aging.
Besides slow relaxation, other glassy properties observed in many circumstances and systems are memory effects and aging, which we now analyze for our systems. As explained in the SI sample excitation, in our case the aging protocol consists in illuminating the sample during a certain period of time t w (called waiting time in the literature) and then switching off the light, letting the sample relax. In this protocol a light intensity smallerthan for the study of the relaxation behavior is used. In Fig. 5A the evolution of the SP is presented for three different values of waiting time t w . The vertical solid line indicates the moment at which the light is switched on, while the dashed lines correspond to the moment when the light is switched off for each t w . It can be seen that during the excitation, the three SP curves approximately overlap. The horizontal solid line shows 〈 S〉 eq , the initial SP value of the sample before any illumination which is used as reference value for the relaxed state.
To study the dependence of the SP relaxation with t w , the data are represented on a logarithmic scale as a function of t/t w taken the time origin (t = 0) as the moment when the light is switched off (Fig. 5B). It is clear that to a good approximation, the three curves collapse. In fact, the overlap curve can be fitted to the equation.
where 〈 S〉 eq is the average SP before illumination, and V 0 is some characteristic excitation surface potential, equals 67 mV in our case. This expression was obtained in the context of relaxation in electron glasses 39,41 , with the factor α = 1. Factors α different from one, ranging between 1 and 10, have been previously observed in electron glasses. For example, in the so-called F protocol, where the sample is driven out of equilibrium by applying a strong non-ohmic electric field 42 . The understanding in the electron glass community is that when the "waiting" period corresponds to relaxation to a new equilibrium state, typically created by changing the gate voltage, one has α = 1, while when the "waiting" period corresponds to a truly excitation phase, as is the case here and in the F protocol, values of α larger than one are found. In Fig. 5A we can also observe that the behavior of the SP when the light is on, the excitation phase, is roughly symmetric to the relaxation behavior, as with the previous relaxation protocol, and a trend often seen in electron glasses 10,14,36 .
It should be highlighted that the present aging protocol leads to important memory effects that prevent a high quality logarithmic relaxation behavior, and therefore it cannot be used for larger t w . In Fig. 5A, one can appreciate in the sample that has been excited for a longer period of time (blue symbols) a small plateau in the SP. This plateau is associated with an incipient memory effect of the excited state. To achieve high quality logarithmic relaxation curves, the "two pass" protocol shown in the previous subsection and higher light intensities should be used. If higher light intensities are used with the present protocol, scaling with t w is still possible, but full aging is lost and one observes subaging, i.e., a scaling in terms of the variable / γ t t w , with the exponent γ smaller than one.
Conclusions
Our SFM studies on the MEH-PPV conducting polymer film show strong evidence of the formation of an electron glass on the surface of the material. This evidence includes the presence of domains on the SP, uncorrelated with topography, and showing self-repulsion, indicative of the relevance of interactions. The fluctuations of SP are compatible with variations of the Coulomb energy of a single charge over the distance between domains. At the same time, the fact that the SP fluctuations are larger than kT and that time correlations are dominated by a broad distribution of characteristic times can be naturally explained within the electron glass model 30,39 .
We have studied the domain structure and dynamics in low degraded (less conducting than non-degraded MEH-PPV). The overall tendency of domain sizes and SP fluctuations is the expected one from our model. The lower the conductance, the smaller the domain size and the larger the variation of the SP. Quantitative comparison between domain size and SP fluctuations in a large enough range of conductance, is outside the scope of the present work and will be addressed in a future work.
The domain dynamics in degraded and non-degraded samples shows similar a similar overal behaviour but a slower dynamics. In addition, it has been previously reported that highly-insulator polymers also present surface charge domains but with no dynamics 43 . These results are consistent with the interpretation that we are observing hopping phenomena.
In addition, after excitation with light, the samples show very slow relaxation of the average SP. Under appropriate excitation conditions, logarithmic relaxation, characteristic of electron glasses, is observed over four decades of time. Relaxation rate decreases with light intensity and so with carrier density. Full aging is also observed: relaxation curves for different excitation times t w can be overlapped when plotted versus t/t w , and the overall relaxation curve follows the prediction for electron glasses, Eq. (2).
Finally, we expect the formation of surface charge domains in hopping-conduction materials to be a very general phenomena. Up to now electron glasses have been studied experimentally through macroscopic properties of the material, in particular conductivity relaxation. We have shown here that the SKPM technique is a powerful tool to monitor and study fundamental properties of disordered systems at the nanoscale and, possibly, at the single charge level. The relation between macroscopic conductivity and the microscopic behavior of slow relaxation is a difficult, but very important issue. In the future we will combine classical macroscopic techniques, mainly based on the analysis of conductivity, with nanoscale techniques to address this problem and to discern between different models for slow relaxation.
Materials and Methods
Sample preparation. Poly SFM Measurements. SFM experiments were performed at room temperature and ambient conditions. For non-optical experiments a Nanotec Electronica SFM setup with a control unit equipped with a PLL/dynamic measurement board was employed. Signal to noise ratio is a limiting factor in our experiments, since high spatial resolution, high energy resolution and large bandwidth are required. In order to have thermal noise-limited performance 44,45 , a monomode polarization preserving fiber is used as light source for the laser detection system. Topography images were acquired working in non-contact dynamic mode using the frequency shift as feedback parameter (FM-DSFM). Frequency modulation SKPM (FM-SKPM) with an AC modulation bias of U ac = 500 mV at 7 kHz was used to measure the local surface potential (SP). Further details of the SKPM set up Scientific RepoRts | 6:21647 | DOI: 10.1038/srep21647 and SFM working modes are described elsewhere 45 . Platinum coated silicon tips (Budget Sensors), with nominal force constant of 3 N/m and resonance frequency of 75 kHz, were used.
Data processing. For the study of equilibrium properties, successive SFM images -movies-have been acquired with constant acquisition parameters. Precise alignment of all movie frames is fundamental, since the processing algorithms need to compare a particular (x, y) position in one frame, with the same (physical) position in other frames. Thermal drift and the length of the movies (up to days) imply that the mechanical setup varies over time and different frames are misaligned. To correct this misalignment and process the images the free WSxM software is used 46 . This software computes the cross-correlation between images and finds the position of the maximum of the cross-correlation corresponding to the offset between images of the movie. This offset defines a drift path which allows computing a drift-free movie.
From these drift-free movies, all relevant information is calculated: single cross-correlation between frames taken at different times, cross-correlation movies obtained by calculating the cross-correlation of each individual frame with a selected frame of the movie (usually the first) and finally average frames. These average frames are calculated by taking the average at each point of any quantity through the whole movie. In particular, the average of the SP is: The well-defined averages of topography and error signal (See Fig. S2 Suplementary information) prove that the drift correction algorithm aligns the frames of a movie to within a few image points; an accuracy of 2-3 pixels is estimated, which is of the order than the spatial resolution of our experiments. We note that exactly the same alignment protocol is applied to all SFM channels.
To analyze dynamical properties, we introduce the time correlation function C(t), defined as ; 5 x y t Sample excitation. Experiments involving sample excitation with light have been carried out in a homebuild SFM implemented in an inverted optical microscope (Nikon Eclipse). As explain elsewhere, this set up allows to illuminate the sample in a controlled way and to measure the topography and SP in darkness as well as under illumination 47 . In the present work, the polymer is illuminated through the transparent ITO:PET electrode with green light (λ = 535 nm) at an intensity between 2 and 7 × 10 16 photons/(s cm 2 ). This wavelength lies within the absorption band of the MEH-PPV inducing exciton generation without degrading the polymer. Under these conditions, only photo-physical reversible processes take place 34,48 . The illuminated area (≈ 300 μm 2 ) is about two orders of magnitude larger than the SFM image size (1-16 μm 2 ). Sample light-excitation experiments have been performed in three steps. First several SKPM images are acquired in darkness for a long period of time (from several hours to days) on a sample region that has been never illuminated before. Secondly, the polymer is illuminated for a certain period of time. Finally the light is completely switched off and the system is left to relax to the equilibrium. During the whole process, the local SP evolution is monitored to study the time evolution.
Two different protocols have been used for sample illumination. In relaxation experiments, the line by line "two pass method" has been used. The detailed working principle and data processing are explained elsewhere 34 . Briefly, in this method the light is pulsed in short "on/off " cycles by scanning each line twice. The first trace is performed in darkness, while the second pass under is illumination. Then the tip moves to the next horizontal line and the process is repeated. In this way, two SKPM images are recorded simultaneously, one in darkness ("off " SKPM image) and the other one under illumination ("on" SKPM image). For the aging experiments a continuous illumination protocol has been used. That is, the light is switched on and the sample is illuminated for a certain period of time (t w ).
To obtain the SP evolution the SKPM images are processed as follows: the mean SP of each horizontal image line is computed and plotted as a function of time. Sample excitation experiments can be performed at a fixed sample position (x 0 , y 0 ), at a fixed sample line (x, y 0 ) or in a sample region (x, y). In this work we have used the latter two cases. It should be noticed that in these cases the SKPM images include not only time dependence information but also spatial information. Therefore the spatially averaged nanoscale behaviour of the SP as a function of time is obtained. | 6,897.8 | 2016-02-25T00:00:00.000 | [
"Physics",
"Materials Science"
] |
An Efficient GUI-Based Clustering Software for Simulation and Bayesian Cluster Analysis of Single-Molecule Localization Microscopy Data
Ligand binding of membrane proteins triggers many important cellular signaling events by the lateral aggregation of ligand-bound and other membrane proteins in the plane of the plasma membrane. This local clustering can lead to the co-enrichment of molecules that create an intracellular signal or bring sufficient amounts of activity together to shift an existing equilibrium towards the execution of a signaling event. In this way, clustering can serve as a cellular switch. The underlying uneven distribution and local enrichment of the signaling cluster’s constituting membrane proteins can be used as a functional readout. This information is obtained by combining single-molecule fluorescence microscopy with cluster algorithms that can reliably and reproducibly distinguish clusters from fluctuations in the background noise to generate quantitative data on this complex process. Cluster analysis of single-molecule fluorescence microscopy data has emerged as a proliferative field, and several algorithms and software solutions have been put forward. However, in most cases, such cluster algorithms require multiple analysis parameters to be defined by the user, which may lead to biased results. Furthermore, most cluster algorithms neglect the individual localization precision connected to every localized molecule, leading to imprecise results. Bayesian cluster analysis has been put forward to overcome these problems, but so far, it has entailed high computational cost, increasing runtime drastically. Finally, most software is challenging to use as they require advanced technical knowledge to operate. Here we combined three advanced cluster algorithms with the Bayesian approach and parallelization in a user-friendly GUI and achieved up to an order of magnitude faster processing than for previous approaches. Our work will simplify access to a well-controlled analysis of clustering data generated by SMLM and significantly accelerate data processing. The inclusion of a simulation mode aids in the design of well-controlled experimental assays.
INTRODUCTION
Cells rely on transmembrane signaling to interact with the outside world. It is essential that cells can specifically and decisively be put into action in response to signals in a noisy and complex environment (Pierce et al., 2002). To do so, mechanisms have evolved that allow the triggering of an all-or-none, lasting response if required. This often involves a threshold number of ligand-activated membrane molecules that recruit auxiliary molecules to form a larger assembly that, upon reaching threshold size, will switch the cell into a different state. These signaling assemblies appear as clusters of membrane proteins in the plasma membrane of cells. However, the clusters may represent only a small subfraction of the membrane protein in question in an otherwise randomly distributed larger population (Janeway et al., 2001;Schultz and Schaefer, 2008). Cluster algorithms can detect such active signaling clusters in a randomly distributed background if the exact spatial distribution of membrane proteins is known (Williamson et al., 2011;Khater et al., 2020). Cartography of membrane protein distribution at the nanoscale has been made possible by super-resolution microscopy approaches based on the sequential localization of single fluorescence-labeled proteins [Single-Molecule Localisation Microscopy (SMLM), Betzig et al., 2006;Rust et al., 2006;Heilemann et al., 2008]. Clustering has since developed into an essential readout for membrane protein function in many cellular processes. Over the last years, several cluster algorithms have been adapted specifically for the analysis of single-molecule fluorescence data of membrane proteins (Owen et al., 2010;Annibale et al., 2011a,b;Nicovich et al., 2017;Baumgart et al., 2019;Arnold et al., 2020;Pike et al., 2020). SMLM of membrane proteins and their cluster analysis still requires a high level of experimental and analytical expertise. To make cluster analysis more accessible, we here combined a selection of the latest clustering approaches with several useful computational features to speed up and streamline cluster analysis in a single, user-friendly software. Specifically, we implemented Bayesian Cluster Analysis, Ripley's-K-based clustering, DBSCAN (Rubin-Delanchy et al., 2015;Griffié et al., 2016), and ToMATo (Pike et al., 2020) for cluster analysis. We then compared the performance of these approaches on simulated and newly generated experimental data from different cellular systems. Furthermore, we implemented a pipeline for parallelized computing of cluster analysis and, as a result, could analyze even large datasets at a fraction of the time required before. Our software will simplify and accelerate cluster analysis as a readout of membrane protein function.
Structure of the GUI
To facilitate the use of parallelized Bayesian cluster analysis for the community, we developed an easy-to-use software called BaClAva (Bayesian Cluster Analysis and visualization application) with a graphical user interface (GUI, Figure 1). This software consists of a pipeline of three modules for simulations, clustering, and analysis that can be used independently via the GUI. Thought experiments are an essential tool in developing reliable experimental strategies and are especially important for data processing-intensive assays because they might offer crucial insights into the experimental setup and data processing strategies. To allow for the freehand design of ground-truth data while simulating realistic experimental output, we included a simulation module similar to FluoSim (Lagardère et al., 2020). This module allows the generation of user-defined clusters of molecules combined with a selected level of randomly placed background molecules. The results of this ground truth are then modeled as images resulting from an SMLM-experiment emulated based on experimental statistics of dye blinking, camera noise, and localization accuracy. The resulting image stack is localized using standard algorithms and can be used as an alternative to or alongside actual SMLM localization data in downstream clustering analysis. If desired, the generation of emulated microscopy images from the constructed localizations can be omitted, as exemplified in Figure 3. This option is based on Griffié et al. (2016).
The second module is the clustering module, which analyzes single-molecule localization datasets in the format [X (nm), Y (nm), STDEV (nm)]. STDEV is the localization precision as calculated by the localization software. Once the data are loaded into the software, the user can choose between ToMATo, Ripley's-K-based, or DBSCAN cluster analysis, define the desired parameter space for Bayesian analysis and select, whether the computation is done sequentially or in parallel. The third and final module allows the visualization and export of the results in a graphic or tabular form, including essential analytical parameters such as the number of clusters, cluster area, and cluster density.
To decrease the number of files stored on the computer disk, we decided to store all information in a Hierarchical Data Format (hdf5) (Figure 1). The hdf5 format enables us to store the localization table (simulation or experimental), the Bayesian engine scores and labels, and further information in a single data file.
Benchmarking
First, we aimed to benchmark our cluster software on simulated clustering data. To do so, we generated 100 simulated images of clustered molecules, each containing ten clusters of 100 localizations. For example, see Figure 2A. These simulations were generated in the following way: Clusters were generated from single points ≥100 nm apart for each of which 100 localizations were generated by drawing from a normal distribution with a standard deviation of 50 nm. The random background was generated at a density of 111 localizations per µm 2 . Thus, the proportion of unclustered localization was designed to be 50% of all localizations (Section 4.6).
These data were then analyzed with the Bayesian model and the three different cluster detection algorithms. Figure 2 shows the simulated data and the corresponding clustering outputs. Since the cluster centers were set to be at least two standard deviations apart from each other, the individual clusters can be correctly identified by eye ( Figure 2A) and as well with DBSCAN ( Figure 2C) and ToMATo ( Figure 2D). In contrast and as shown before (Pike et al., 2020), the approach based on Ripley's K-function ( Figure 2B) fails to separate nearby clusters and thus commonly misidentifies cluster number and area ( Figures 2E,F). As previously shown, this behavior is due to the incapability of this approach to correctly take into account the local density of the data points (Rubin-Delanchy et al., 2015;Griffié et al., 2017). In contrast, both DBSCAN and ToMATo could quantify both cluster number and overall cluster area quite accurately in the majority of simulations ( Figures 2E,F). These methods in the Bayesian cluster approach rely not on a single set of parameters but instead on a continuum of so-called proposals, defined sets of values computed to cover an ample parameter space to find an overall optimum of cluster identification (Rubin-Delanchy et al., 2015;Griffié et al., 2016). While this approach has proven to lead to superior results, it is necessarily computationally costly. We aimed to overcome this problem to increase processing speed and thus experimental throughput.
In the original work (Rubin-Delanchy et al., 2015;Griffié et al., 2016), the cluster proposals' calculation in Bayesian analysis is done in nested for-loops on a single CPU core. Since the individual cluster proposals are independent of each other, the processing could also be implemented in parallel. This means that the program uses multiple CPU cores instead of a single core and therefore calculates multiple proposals at the same time. In our software, we implemented the parallelized computing of Bayesian cluster analysis and compared the results with the sequential computational approach.
We first used ten simulations to benchmark the clustering methods described above in Bayesian analysis. We found that typical runtimes for Ripley's K-based and DBSCAN clustering were 25.78 ± 0.86 and 28.45 ± 0.78 min, respectively (mean ± standard deviation). The ToMATo implementation from the RSMLM package (Pike et al., 2020) had a runtime of 23.87 ± 0.80 min (mean ± standard deviation, Figure 3). By parallelizing the clustering and scoring process to multiple cores, we found the computation time to decrease by 60% for Ripley's K-based, 10.41 ± 0.23 min, and DBSCAN, 11.90 ± 0.27 min (Figure 3). For the ToMATo implementation, the computational time decreased by one order of magnitude to 3.062 ± 0.072 min. In summary, the parallelization significantly reduced processing time for Bayesian cluster analysis.
FIGURE 1 | Overview of the software GUI. Schematic of the three independently usable software modes and organization of the software. Simulations can be prepared individually or as batches, and the localization results get exported as tiff or hdf5 files, depending on the simulation option. For the second module, the simulated data is imported from the hdf5 file, or experimental datasets can be imported in the form of a text or csv file. The user can set various parameters, most notably the cluster method, the type of computation, and additional Bayesian clustering parameters. This module's output, namely the scores and the labels for all proposals, is stored in the hdf5 file. In the final processing step, the original localization table and the Bayesian clustering module's output are used to produce the best cluster plots and the corresponding (batch) statistics. The statistics are exported as text files as well as plots.
Frontiers in Bioinformatics | www.frontiersin.org October 2021 | Volume 1 | Article 723915 The mean is emphasized as a black circle. Ten clusters were simulated, and the mean for Ripley's-K-based clustering was 9.8 ± 2.0, for DBSCAN 9.5 ± 0.7, and 9.8 ± 0.7 for ToMATo. Note that the spread is significantly larger for Ripley's-K-based, DBSCAN never overcounted, and ToMATo was the most accurate overall. (F) Plot of all ground truth and recognized cluster areas. The ground truth data's cluster area has an average size of 0.061 ± 0.013 µm 2 , the Ripley's-K-based clustering results in 0.044 ± 0.023 µm 2 , DBSCAN in 0.055 ± 0.017 µm 2 , and ToMATo clustering averages the area to 0.053 ± 0.015 µm 2 (mean ± standard deviation). Next, we aimed to investigate several known sources of error in clustering single-molecule localization microscopy data. An important source of error in the cluster analysis of SMLM data is caused by multiple localizations of the same fluorescent molecule generated by most SMLM approaches that necessarily generate a cluster of localizations from every single fluorophore. Consequently, this fact must be considered for declaring any statement on fundamental information such as cluster size in terms of area and number of molecules in the cluster. In the case of PALM, algorithms have been published which aim at correcting this artifact (Annibale et al., 2011b,a;Jensen et al., 2021b,a). By simulating blinking SMLM data with realistic blinking statistics for Alexa Fluor 647, we determined how dense the underlying molecules must be for proper cluster detection. The simulations of (d)STORM experiments were generated in the following way: Cluster areas were generated by randomly distributing 40 non-overlapping clusters with an area of 0.0078 µm 2 (diameter 50 nm). Their molecular density was increased from 0.71 ± 0.25 × 10 3 to 6.24 ± 0.63 × 10 3 μm −2 , translating to molecules per cluster ranging from 5.6 ± 1.9 up to 49.0 ± 4.9. The random background was generated at a density of 639 ± 49 molecules per µm 2 for sparse clusters up to a density of FIGURE 4 | Influence of fluorophore blinking on clustering. (A) Violin plot for the relative density of the clusters vs. the background with and without grouping applied, (B) Violin plot of the percentage of the clustered localization with grouping, (C) Violin plot of the number of clusters per ROI with and without grouping applied, (D) Violin plot of the areas of the clusters with and without grouping, (E) Examples for clustering of a random distribution of fluorophores and 40 clusters at a density of 1.40 ± 0.36 × 10 3 μm −2 (left column), 3.79 ± 0.38 × 10 3 μm −2 (middle column) and 6.24 ± 0.63 × 10 3 μm −2 (right column). Each dataset was analyzed in SMAP with and without grouping. The cluster analysis was performed by the Bayesian engine plus ToMATo.
Frontiers in Bioinformatics | www.frontiersin.org October 2021 | Volume 1 | Article 723915 346 ± 51 molecules per µm 2 for dense clusters. For sparse clusters, almost 94 and 51% of all molecules are assigned to the background for dense clusters. The blinking parameters were k on 0.01 s −1 and k off 10 s −1 . The FWHM of the PSF was set to 200 nm with an intensity of 2007. The pixel size of the camera was set to 0.096 µm, which is identical to the pixel size of the Evolve Delta 512 Photometrics camera on our microscope. The exposure time was set to 10 ms, which is the exposure time we use in experiments with Alexa Fluor 647 dyes, and as in a (d)STORM experiment, 50,000 frames were acquired. The localization procedure and grouping were done in SMAP (Ries, 2020). The obtained localization table was used for the Bayesian Analysis.
The results are visualized in Figure 4.
In Figure 4A, the cluster to background density for grouped and non-grouped data is shown. For both cases, the relative density increases with increasing cluster density and a smaller spread of the distributions for grouped data, whereas the nongrouped data distributions show a broader spread, indicating the efficiency of the grouping function in SMAP. Additionally, the number of clusters per region of interest (ROI) is reduced by the grouping, which removes clusters caused by single blinking fluorophores ( Figures 4C,E). The higher relative density of fluorophores in clusters compared with background localizations indicates that a local density threshold must be surpassed to render the interpretation of cluster data independent of fluorophore blinking properties. As shown in Figure 4C, the number of clusters is constant for grouped data up to a concentration of 2.62 ± 0.39 × 10 3 μm −2 localizations. For higher concentrations, the number of clusters approaches the ground truth of 40 clusters. Without grouping, the number of identified clusters decreases with increasing fluorophore concentration, reflecting a higher relative enrichment of fluorophores inside the clusters than outside them. The improved situation for grouped data is also visible in Figure 4B, showing that the percentage of clustered localizations increases with increasing fluorophore density. For the best cluster result in these simulations, more than 30% of the localizations must occur in clusters, and a relative density (localization density inside vs. outside of cluster) threshold of 10 must be overcome for the localizations inside clusters versus outside.
Moreover, the cluster size ( Figure 4D), meaning the area covered by localizations in a cluster, shows the influence of background localizations on the data distribution. Cluster area increases in size for grouped data starting from a concentration of 1.87 ± 0.39 × 10 3 μm −2 . For the non-grouped data, there is a significant proportion of very small clusters at all concentrations. This cluster population is not present for the grouped data, indicating that these clusters emerge from multiple detections of a single fluorophore, i.e., blinking. For a density of 2.62 ± 0.39 × 10 3 μm −2 molecules and higher, a second population emerges in the non-grouped data distributions, which corresponds to the main population in the grouped distributions. Therefore, they can be considered correctly identified clusters. Similarly, from 2.62 ± 0.39 × 10 3 μm −2 molecules onwards, the number of clusters per ROI decreases. As demonstrated in Figure 4E, small background clusters are removed with the grouping functionality (top row vs. bottom row) and with increasing fluorophore density within the clusters (from left to right). As expected, the ground truth clusters become more apparent when the number of clustered molecules is increased even in the non-grouped data, indicating that single fluorophore blinking has a significantly reduced impact on density-based cluster identification for denser clusters. We concluded that grouping is essential in the detection of smaller clusters.
Finally, we aimed to apply our algorithm to experimental data from single-molecule localization experiments of intact cells. We used standard controls in the field for non-clustered and clustered molecules respectively at the plasma membrane. The lipidanchored glycosylphosphatidylinositol-coupled green fluorescent protein GPI-GFP should be more or less homogeneously distributed and functioned as the negative control. The clathrin-light chain (CLC), of which dozens of copies are incorporated into every ∼150 nm diameter clathrincoated pit and thus appears strongly clustered, served as the positive control. In order to keep our results comparable, all molecules of interest were tagged with a GFP protein, and the (d) STORM dye Alexa Fluor 647 was bound to the GFP via anti-GFP nanobodies in all experiments (Ries et al., 2012). From the simulation work, we know that the cluster results for GPI-GFP should show a wide range of cluster areas, whereas, for the CLC, we expect to yield well-defined cluster areas. Finally, we asked whether we could detect clustering for the transmembrane receptor CD95, as the receptor activation via its ligand may trigger apoptosis or tumorigenesis of cancer cells and has been suggested to result in the formation of high order molecular clustering (Martin-Villalba et al., 2013). CD95 was likewise labeled via GFP and AF647 nanobodies.
The reconstructed images in Figure 5 of these three proteins show differences in the spatial distribution of the localizations. For GPI-GFP imaged in CV-1 cells in Figure 5A, the localizations are evenly distributed, and the cluster maps for the zoom-ins show small clusters, which are probably due to the blinking of the Alexa Fluor 647 dye. In contrast, in Figure 5B, the CLC imaged in HeLa cells show well-defined clusters in agreement with clathrincoated pit size (Supplementary Figure S1) with little background localizations, as seen in the cluster maps of the zoom-ins. The CD95 receptor in T98G glioblastoma cells presents a localization distribution with smaller clusters and more background localizations than CLC. The cumulative distribution of the cluster areas of several cells for each condition in Figure 5D reveals that GPI and CLC exhibit distinct distributions of their respective cluster areas in agreement with expectations. The cumulative distribution of cluster areas for the CD95 receptor is positioned between the two controls, demonstrating that CD95 forms small clusters likely consisting of around 0.54 molecules/ nm 2 in the plane of the membrane.
DISCUSSION
Here we present a user-friendly software solution for cluster analysis of SMLM data. Our software significantly reduces processing time and allows the user to select different algorithms to identify and quantify cluster formation. The simplest cluster algorithms, such as nearest-neighbor algorithms, can answer whether areas of above-average concentration, clusters, exist in a field of view (Endesfelder et al., 2014). For a more detailed analysis of clusters found in cellular membranes, Ripley's K-function can provide answers to the length scale of interparticle distances and the proportion of entities found in clusters in a given dataset (Owen et al., 2010). However, these methods are prone to artifacts intrinsic to singlemolecule fluorescence-based microscopy approaches, which lead to small local clusters due to the blinking behavior of individual fluorophores. To overcome errors due to blinking, approaches have been developed to determine the degree of clustering in challenging experimental circumstances, such as for dense membrane molecules by varying the dye density (Baumgart et al., 2019;Arnold et al., 2020).
To understand the functional underpinnings of cluster formation in cell biology, a qualitative view of clustering is not sufficient, but reproducible, robust quantitative assays are required. One of the first ideas put forward was to use Ripley's function not on the entire sample but on individual localizations convoluted with a search radius and a clustering threshold for cluster identification in dense backgrounds. To facilitate the differentiation of Ripley's K function on the entire sample or individual localizations, we termed the second case Ripley's-Kbased clustering. Ripley's K-based clustering cannot adequately determine clusters in samples with variations in cluster density and final cluster size. Therefore, density-based clustering methods, like DBSCAN (density-based spatial clustering of applications with noise), have been adapted for SMLM cluster analysis, and they are less error-prone, as shown before (Pike et al., 2020). In DBSCAN, the search parameters are the search radius and the minimal number of data points within this search radius. Points for which these parameters are valid are counted towards this cluster. After identifying all clusters, the remaining data points are assigned to the background, and other parameters such as the individual cluster area and density can be extracted from the data.
Even though density-based methods, like DBSCAN, can handle datasets with density variations, they still often fail to separate individual clusters that are very close to one another but are easily distinguishable by eye. Additionally, it has been shown that the detection of clusters by DBSCAN and Ripley's-K-based clustering (Pike et al., 2020) can be sensitive to even small changes in the analytical parameters, possibly leading to artifactual results. One of the latest introductions to SMLM cluster analysis is persistence-based clustering which is based on density estimation . The introduced scheme is called Topological Model Analysis Tool (ToMATo, Pike et al., 2020;Chazal et al., 2012Chazal et al., , 2013. In contrast to the abovementioned density-based methods, in this algorithm, local maxima in molecular density are identified and termed clusters by introducing a density gradient generated by creating a path emanating from a molecule to its neighbors and using the intermolecular distance as a measure of density. If this intermolecular distance increases, the border of a cluster may be reached. Such a local maximum in distance or minimum in density may be a saddle point between clusters or define the outer perimeter of a cluster. A threshold value defines the persistence of a cluster from its center into space. Clusters with persistence smaller than the threshold are assigned to neighboring clusters, or they are deemed background. As a result, ToMATo allows for a separation of even partially overlapping clusters, and additionally, the output clustering results are less sensitive to analytical input parameters as compared to Ripley's K-based clustering and DBSCAN. Further clustering methods for SMLM data are based on Voronoï tessellation (Levet et al., 2015), which detects clusters based on polygonal regions. Voronoï tessellation intrinsically generates contours of regions of density that may also be used for boundary detection of cells. Recently, machine-learning has been employed to improve cluster detection (Williamson et al., 2020), but the number of input neurons limits the correct processing of the underlying information. The Bayesian engine's main drawback is that due to the calculation and scoring of thousands of cluster proposals for optimal results, the process is significantly slowed down compared to traditional methods with a single cluster proposal, hampering the routine use of this method. On the other hand, ToMATo clustering parameters are determined based on a persistence diagram which can cause user bias.
To overcome these limitations and provide accessible GUIbased software for state-of-the-art cluster analysis, we implemented Ripleys-K-based clustering, DBSCAN, and ToMaTo in a common software that allowed for parallel computing. In this software, we first improved the Bayesian engine's speed by implementing parallel computation and introduced ToMATo clustering to the Bayesian engine, thereby dramatically decreasing computational time. In combination with the software GUI, the Bayesian engine has an improved user experience and processing speed, which we hope will make state-of-the-art SMLM cluster analysis available in many laboratories.
During an SMLM measurement of several thousands of frames, a fluorescent molecule may cycle several times between a bright and a dark state, and thus, one molecule may be detected multiple times within a radius determined by its localization precision. As a result, it is impossible to differentiate between a single molecule detected several times in different frames and different molecules in close proximity detected in different frames. This is, of course, especially problematic in cluster analysis, where localizations are processed first without bias. To overcome this problem, it is important to develop an understanding of the degree of influence of blinking in the dataset at hand. As we showed in Figure 4, the number of the small clusters resulting from dye blinking decreases with increasing molecular density within the clusters while keeping the actual cluster size constant. Thus, there is an intrinsic threshold for relative localization densities inside and outside clusters that render blinking irrelevant. This holds true under the assumption that all localizations are caused by dyes bound to molecules of interest, and no false localizations are present in the sample. Below this threshold, the number of detected clusters is highly overestimated, and the cluster radii are dramatically underestimated. From simulations, we know that such single-molecule clusters can be detected as sub-peaks within clusters at low-density ratios. Increasing the density ratio now increases the chance that clusters are quantified at their true size. It is common in SMLM data analysis that multiple temporally and spatially closely correlated localizations are grouped together in a final reconstruction and are thus counted as a single molecule. In clustering, this procedure reduces the number of small background clusters dramatically, and we analyze this effect in depth in Figure 4. Our grouping is based on the blinking behavior of the most used (d)STORM dye Alexa Fluor 647, which we also used in our experiments. Likewise, we also based our simulations on the blinking behavior of this dye (Heilemann et al., 2008;Dempsey et al., 2011). In order to detect clusters of smaller density ratios and smaller sizes either or both, the cluster detection may be improved by changing the dye or even the SMLM method from (d)STORM to DNA-PAINT as shown in Jayasinghe et al. (2018).
Microscopy experiments in cells are much more complex than the corresponding in-silico experiments because many different known and unknown cellular processes are involved in the Frontiers in Bioinformatics | www.frontiersin.org October 2021 | Volume 1 | Article 723915 temporal and spatial organization of the cell molecules and may interfere in the process studied. Therefore, we chose highly abundant molecules as cellular controls for clustering experiments. A simple positive control is clathrin-coated pits expressed at a well-defined radius of around 80 nm (A ≈ 0.02 µm 2 ) (Sochacki et al., 2017). Negative controls for clustering in a cellular environment are far more challenging to identify because natural cellular signaling processes result in a spatial and temporal reorganization of the involved molecules, and many membrane molecules exhibit clustering of some extent (Gowrishankar et al., 2012;Saka et al., 2014;Baumgart et al., 2016;Kalappurakkal et al., 2020). Therefore, the influence for altering the negative control's organization by cell processes should be kept at a minimum, and an artificially introduced protein that is only anchored to the outer membrane of the plasma membrane and has no natural interaction partners, such as GPI, is the ideal option (Li et al., 2020). These extreme cases of clustering and nonclustering probes can be well differentiated in their reconstructed images as well as their cumulative distribution functions. Proteins with so far unknown spatial distribution on the plasma membrane, such as the transmembrane receptor CD95, should present a behavior between these two extremes. If they are less clustered, they should tend towards a behavior similar to GPI, and with increasing cluster areas, they should tend towards a distribution similar to clathrin-coated pits. Since CD95 can be found at the plasma membrane as monomers or homodimers and homotrimers (Micheau et al., 2020), it should be detected as small clusters, as observed in Figure 5. We conclude that our software can correctly distinguish between unclustered molecules and clusters of even small size and a few molecules in number. Taken together, our work allows the implementation of singlemolecule clustering analysis at a high rate of data throughput for beginning users. We expect our work to accelerate research in this area significantly and to contribute to the acceptance of reproducible standards in clustering data analysis. In future work, other analytical methods such as Voronoï tessellation (Andronov et al., 2018;Levet et al., 2015) and extensions to 3D (Griffié et al., 2017) and dual-color co-clustering (Jayasinghe et al., 2018) may be implemented, and the processing speed may be further improved, i.e., by the implementation of GPUprocessing.
Cell Culture and Preparation
CV-1 cells were cultured in a standard DMEM medium (1X, Gibco) supplemented with 10% FBS (ThermoFisher) and 1% GlutaMax (100X, Gibco by Life Technologies). Stable HeLa CLC-GFP cells were cultured in the same medium with an additional 1% Penicillin-Streptomycin (Sigma), and for the T98G CD95-GFP cells, 1% sodium pyruvate (stock: 100 mM, Gibco) was added to the medium. The vector CD95-GFP was infected into the cells with a lentiviral construct. The cells were then FACS sorted for the stably transfected clones. All cell lines were regularly tested for mycoplasmas and only used when tested negative. For the seeding of the cells, 18 mm diameter #1.5 glass slides (VWR) were cleaned in an ultrasound bath for 20 min using 2% Hellmanex III (Hellma) and 70% ethanol, respectively. Afterward, the glasses were dried and plasma cleaned for another 30 min.
CV-1 GPI-GFP Cells
Transfection of GPI-GFP into CV-1 cells was done with lipofectamine 3000 following the standard protocol (lipofectamine protocol by Invitrogen/ThermoFischer). Cells were treated with trypsin-EDTA and seeded on the glass slides for incubation of 24 h (densities: 6 × 10 6 cells/ml for CV-1, 7 × 10 4 cells/ml for HeLa CLC-GFP and T98G CD95-GFP). The transfected CV-1 cells were fixed with prewarmed 4% PFA with 0.2% GA in PBS for 20 min at 37°C. Then, cells were quenched with freshly prepared 0.1% NaBH 4 in PBS for 7 min at room temperature and extensively washed. Cells were blocked in two steps: for 30 min with ImageIT, followed by 4% goat serum in 1% BSA in PBS for 1 h. CV-1 GPI-GFP cells were stained with anti-GFP nanobodies (FluoTag-Q anti-GFP) labeled 1:1 with Alexa Fluor 647 from NanoTag Biotechnologies GmbH at a concentration of 50 nM for 1 h. Afterward, cells were postfixed with 4% PFA and 0.2% GA in PBS for 20 min and quenched with 0.1% NaBH 4 in PBS for 5 min at room temperature.
HeLa CLC-GFP Cells
HeLa CLC-GFP cells were fixed with prewarmed 4% PFA in PEM for 20 min at 37°C and quenched with NH 4 Cl in PBS for 5 min at room temperature. After quenching for 5 min with 0.2% saponin in PEM, the cells were blocked with 4% goat serum in 1% BSA in PEM for 1 h. HeLa CLC-GFP cells were stained with the NanoTag Biotechnologies GmbH nanobody for 30 min at a concentration of 50 nM and afterward post-fixated with 4% PFA in PEM for 20 min at room temperature. The cells were quenched with NH 4 Cl in PBS for 5 min. In between all steps, the HeLa cells were extensively washed with PEM.
T98G CD95-GFP Cells
The T98G CD95-GFP cells were fixed for 20 min at 37°C with prewarmed 4% PFA plus 0.2% GA in PEM and quenched with freshly prepared 0.1% NaBH 4 in PEM for 7 min T98G cells were permeabilized with 0.2% saponin in PEM for 5 min and blocked with 4% goat serum in 1% BSA/PEM for 1 h. The cells were stained with the NanoTag Biotechnologies GmbH nanobody at a concentration of 50 nM for 30 min and post-fixated with 4% PFA with 0.2% GA in PEM for 20 min at room temperature. For the post-quenching, the cells were incubated in 0.1% NaBH 4 in PEM for 7 min. In between all steps, the cells were extensively washed with PEM.
(d)STORM Imaging
The fixed and stained samples were mounted and imaged in betamercaptoethanol and GLOX (2.5 mg/ml glucose oxidase, 0.2 mg/ ml catalase, 200 mM Tris-HCl pH 8.0, 50% glycerol) as imaging buffer (10:1). The (d)STORM images were acquired on a home build TIRF microscope as described in (Albrecht et al., 2016) objective was used to reach a pixel size of 96 nm. The samples were illuminated with a 639 nm laser (Changchun New Industries Optoelectronics Tech. Co., Ltd.) at powers of 0.008-0.015 mW/ μm 2 . For the acquisition of the (d)STORM images, a water-cooled and back-illuminated Photometrics EMCCD camera with 512 × 512 pixels at a pixel size of 16 × 16 µm was used for the acquisition of 30,000 frames at an exposure time of 10 ms. The EMCCD camera was calibrated before the data acquisition, and the image acquisition was controlled with MicroManager.
(d)STROM Reconstruction
The acquired and simulated (d)STORM datasets were localized using SMAP (Ries, 2020). Important camera and acquisition parameters were extracted from the metadata file, which had been saved with the data. Furthermore, the electron multiplier (EM) gain was set to 300, and the conversion factor to 6.7 (analog to digital units to photons). The minimum distance between two candidate peaks in order to be fitted separately was set to 7 pixels. For the point-spread function (PSF) fitting, the following parameters were set to a differential of Gauss with sigma 1.2, dynamic factor 1.7, and free PSF, using the workflow "set Cam parameters."
Grouping in SMAP
The grouping procedure is a part of SMAP, which we used for the reconstruction. The number of frames, dT, for which a single molecule can be non-fluorescent but still be grouped with the first localization of that molecule was set to dT 1. The distance, dX, the centroid of a single molecule can be shifted in the image plane between two consecutive frames, but still, be grouped with the first localization of the molecule, was set to dX 1. These are the standard values in SMAP, and they were identified as the optimal parameter values for out (d)STORM experiments.
Simulations Used for Cluster Algorithm Comparison
The simulations were done with an adapted simulation code published by Rubin-Delanchy et al. (2015). The number of clusters, the number of molecules inside each cluster, the corresponding standard distribution for the cluster size, and the background percentage were set depending on the analysis. Unlike the original publication, the cluster centers are set to be at least two standard distributions apart from each other. In total, 100 simulations were done for each case.
Simulations Used for Computational Time Evaluations
Ten simulations with a standard deviation of 50 nm, 10 clusters with 100 molecules each, and 50% of the total number of localizations in the background were used to determine the computational cost for the three cluster algorithms combined with the Bayesian engine. The field of view had a size of 3,000 × 3,000 nm 2 , and the background is uniformly distributed. The localization precisions are generated from a gamma function with shape 5 and rate 0.166667 (default parameters, Griffié et al., 2016).
Simulations With Blinking Molecules
Simulations were prepared in Fluosim (Lagardère et al., 2020). For the simulation of the sample staining, a geometry file was created with a python script. The field of view had a size of 25 × 25 μm 2 and was composed of 40 randomly distributed, nonoverlapping clusters with a diameter of 50 nm. The clusters were positioned with a minimum distance of 500 nm from any border of the sample. The background image was an image of the Evolve 512 EMCCD camera (Photometrics) with a size of 26 × 26 μm 2 . The pixel size matched the pixel size of our experimental setup. Each pixel's noise values were not considered because only the pixel shape was used in the further course. The number of molecules was set to 4,000 to match the density of optimal CV-1 GPI-GFP samples stained with anti-GFP nanobodies labeled with Alexa Fluor 647. For a fixed period (5-50 s), the molecules were diffusing within the field of view with a coefficient of 0,01 μm 2 /s. A binding rate of 0,997-1,007 s −1 was set to allow cluster formation inside the clusters. Outside the designated cluster areas, the binding rate was set to 0 s −1 . After the binding period, the molecules were freely diffusing for 50 s. During this time, the binding and unbinding rates within the clusters were set to zero and set to 0,997-1,007 s −1 outside of the cluster areas, thereby causing a homogeneous distribution of background molecules.
For simulating an actual SMLM experiment, the fluorophores' blinking parameters and the optical properties of the fluorescence emission were set accordingly. The on-rate was 0.01 s −1 , and the off-rate 10 s −1 , based on an estimated 1:1,000 ratio in an SMLM experiment. For fitting the point-spread function, full-width at half maximum was fixed at 200 nm with a fluorescence emission intensity of 2007. As in a microscopy experiment, 5,000 frames were acquired of the simulated sample, and the exposure time was set to 10 ms/frame. The output tiff-file was localized in SMAP with the standard parameters used for SMLM imaging. The camera parameters were the default values of the Delta 512 as given by its metadata file.
Computational Runtime Measurements
To evaluate the implemented cluster algorithms' speed, we used a standard 64-bit laptop computer running Linux (Ubuntu 18.04.5 LTS), equipped with GNOME 3.28.2, 7.7 GiB of memory, and 4 Intel ® Core ™ i5-6200U CPU @ 2.30 GHz processors. The R library "tictoc" (Izrailev, 2014) was used to measure the time needed for each dataset to be processed.
Cluster Algorithms
The Ripley's-K-based and DBSCAN cluster algorithms used were written by (Rubin-Delanchy et al., 2015;Griffié et al., 2016). The code was adapted for improvement by using functions from several R packages and the ToMATo cluster algorithm for SMLM data adapted from the R package RSMLM (Pike et al., 2020). The library "doParallel" was used for parallel implementation (Analytics and Weston, 2014).
Bayesian Parameters
All Bayesian cluster scorings were done with the same set of parameters. The percentage of background localizations was set to 50%, and the Dirichlet process's concentration coefficient was Frontiers in Bioinformatics | www.frontiersin.org October 2021 | Volume 1 | Article 723915 20. The optimal cluster parameters (radius and threshold) were searched in the sequences 5 to 300 for the first parameter and 5 to 500 for the second parameter in steps of 5.
Statistical Analysis
The statistical comparison was performed with a self-developed R script.
DATA AVAILABILITY STATEMENT
The raw data supporting the conclusions of this article will be made available by the authors, without undue reservation. | 9,266.6 | 2021-06-11T00:00:00.000 | [
"Computer Science",
"Biology"
] |
A New [ 14 8 3 ]-Linear Code From the Aunu Generated [ 7 4 2 ] -Linear Code and the Known [7 4 3 ] Hamming Code Using the (U|U+V) Construction
In this communication, we enumerate the construction of a [ 7 4 2 ]linear code which is an extended code of the [ 6 4 1 ] code and is in one-one correspondence with the known [ 7 4 3 ] Hamming code. Our construction is due to the Carley table for n=7of the generated points of was permutations of the (132) and (123)-avoiding patterns of the non-associative AUNU schemes. Next, [ 7 4 2 ] linear code so constructed is combined with the known Hamming [ 7 4 3 ] code using the ( u|u+v)-construction to obtain a new hybrid and more practical single [14 8 3 ] errorcorrecting code. Citation: Ibrahim AA, Chun PB, Kamoh NM (2017) A New [ 14 8 3 ]-Linear Code From the Aunu Generated [ 7 4 2 ] -Linear Code and the Known [7 4 3 ] Hamming Code Using the (U|U+V) Construction. J Appl Computat Math 7: 379. doi: 10.4172/2168-9679.1000379
Introduction
Historically, Claude Shannon's paper titled "A Mathematical theory of Communication" in the early 1940s signified the beginning of coding theory and the first error-correcting code to arise was the presently known Hamming [7,4,3] code, discovered by Richard Hamming in the late 1940s [1]. As it is central, the main objective in coding theory is to devise methods of encoding and decoding so as to effect the total elimination or minimization of errors that may have occurred during transmission [2] due to disturbances in the channel. The special class of the (132) and (123) avoiding Patterns of AUNU permutations has found applications in various areas of applied Mathematics [3]. The authors had reported the application of the adjacency matrix of Eulerian graphs due to the (132) -avoiding patterns of AUNU numbers in the generation and analysis of some classes of linear and cyclic codes [4,5], respectively. The authors utilized the Carley tables for n=5 [6] to derive a standard form of the generator/parity check matrix for some code. In this article [7], we enumerate the construction of a [7 4 2] -linear Code from the Carley table for n=7 of the generated points of was permutations of the (132) and (123) avoiding patterns of the non-commutative AUNU schemes [8]. The [742] linear code is then shown to be an extended [9] code of the [ 6 4 1 ] code and is in oneone correspondence with the [ 7 4 3 ]-Hamming Code. Moreover, the [ 7 4 2 ]-linear code so generated is [10,11] then combined with the known Hamming [ 7 4 3 ] code using the (u|u+v) construction method to obtain a new and more Practical single error correcting code with dimensions [12] n=14, k=8 and d=3.
Generator matrix
A generator matrix G for a linear code C is a kXn matrix for which the rows are a basis for C. If G is a generator matrix for C, then . G is said to be in standard form(often called the where I k is the k×k identity matrix.
U U V construction
Two codes of the same length can be combined to form a third code twice the length in a way similar to the direct sum of the codes construction. This is achieved as follows; Let C i be an [n,k i ,d i ] code for i∈{1,2}, both over the same finite field
Remark
If C i is a linear code and has generator matrix G i and parity check matrix H i , then the new code C as defined in (1) above has generator and parity check matrices as, 1 1 1 Since the minimum distance of the direct sum of two codes does not exceed the minimum distance of either of the codes, then it is of little use in applications and is primarily of theoretical interest. As such, we for the purpose of this research concentrate on the ( ) u u v + construction.
Methodology Cayley tables
We consider the Carley table below, which is constructed using A n (132) for n=7 [2] We now convert the entries of the Carley table above to the binary system using Modulus 2 arithmetic. Table 1, thus becomes; The above table is the matrix G below; Clearly, all possible linear combinations of the rows of G generates a linear Code say C of length n=7 and size M=16 with the following Code words; Now, since G has five rows and generates a code with sixteen code words, we seek a generator matrix for G. Deleting the first row and last column of G, we obtain a Matrix say 1 0 1 1 0 0 1 1 0 0 1 1 1 0 0 0 1 1 Next, we apply the following series of row operations on G II , we have; Respectively as follows Observe that G I above is a generator matrix in standard form [1], i e G I X × − = , but G I which is a generator matrix for C has words (rows) of length n=6, while C has word length n=7. Note also that all words in C has even weight, therefore extending the rows of G I by adding one digit each, to make the number of nonzero coordinates in each row of G I even, we obtain; Which is the required generator matrix for our [7 4 2]-linear code C.
On combining the two codes using the ( ) u u v + construction, we shall obtain a (123)-avoiding patterns of the non-associative AUNU schemes has been used to construct a [742]-linear code C which is an extended code of the [641] code. Moreover, the [742]-linear code so generated is then utilized using the (u|u+v) construction method to obtain a new and more practical single error correcting code with dimensions n=14, k=8 and d=3.
Conclusion
This paper has further pointed the applicability of the AUNU schemes in the direction of coding theory, i.e., Codes, generation and analysis. | 1,435 | 2018-01-01T00:00:00.000 | [
"Mathematics"
] |
Dual-Pump Approach to Photon-Pair Generation: Demonstration of Enhanced Characterization and Engineering Capabilities
We experimentally study the generation of photon pairs via spontaneous four-wave mixing with two distinct laser pulses. We find that the dual-pump technique enables new capabilities: 1) a new characterization methodology to measure noise contributions, source brightness and photon collection efficiencies directly from raw photon-count measurements; 2) an enhanced ability to generate heralded single photons in a pure quantum state; and 3) the ability to derive upper and lower bounds on heralded-photon quantum state purity from measurements of photon-number statistics even in the presence of noise. Such features are highly valuable in photon-pair sources for quantum applications.
INTRODUCTION
Optical quantum states for quantum applications are commonly realized through photon-pair generation using spontaneous parametric down-conversion (SPDC), in which one pump photon is annihilated and a photon pair is created via the second-order nonlinearity in an interaction medium, and spontaneous four-wave mixing (SFWM), in which two pump photons are annihilated and a pair is created via a χ (3) nonlinearity. Much research has been conducted on engineering and optimizing these techniques for quantum applications such as quantum computation [1], quantum metrology [2], and quantum communication [3]. Important figures of merit include the pair-generation probability and the collection efficiency, which determine bit rates and signal-to-noise ratio, and false detections, which degrade the quality of the photon-pair source. Despite the importance of these parameters, in many sources noise contributions prohibit their direct assessment, and to date such photon-pair source characteristics have been calculated by neglecting noise or assuming pump-independent noise contributions [4], indirect measurements of various source performances [5][6][7], or from fits to multiple measured data points [8][9][10]. Another figure of merit that is critical for protocols and gates that rely on interference of photons from separate sources is the degree of quantum state purity of the individual photons.
Spontaneous four-wave mixing, unlike the SPDC process, can occur with two spectrally distinct pump fields (see Fig. 1(a)); this additional degree of freedom is beneficial for photon-pair source design, with experimental uses including degenerate photon-pair generation [11][12][13][14] and avoiding single-pump SFWM background [15]. In addition, it has been shown theoretically that dualpump SFWM leads to improved capabilities in tailoring the inter-correlations of the photon pairs [16,17], including a proposed method that relies on the group-velocity difference between pump pulses [18].
Here we report experimental demonstrations of some key advances in photon-pair generation in general, and dual-pump SFWM in particular. First, we show how the dual-pump scheme enables a simple and direct measurement of the noise contribution to the detection events; this noise consists of background photons from ambient light, photons from additional processes that occur concurrently with photon-pair generation, or false detection events. In turn, measurement of the noise contribution allows a direct quantitative assessment of source performance, including the photon-pair generation rate as well as overall collection and detection efficiencies of the created photons [5]. Second, we show that the group delay between the two pump pulses enables the creation of photon pairs where each of the individual photons is in a highly pure quantum state [18], and third, we derive a new way to determine the lower and upper bounds for the individual photon purity from second-order coherence measurements in the presence of noise -a method that naturally applies to the dual-pump SFWM where the noise can be directly characterized, but may also find use to estimate other types of photon-pair sources. This Paper is organized in the following way: in Section we give an overview of photon-pair generation in dual-pump SFWM, including generation probability and a description of the quantum state. In Section we demonstrate and analyze photon-pair production with dual pumps. In Section we experimentally confirm the advantage of the dual-pump scheme in generating photon pairs with reduced spectral correlations, both through joint spectral density and single-photon purity measurements. Finally, in Section we conclude and discuss our results.
1 and pump 2 with carrier angular frequencies ω p1 and ω p2 , respectively -enter a χ (3) medium where one photon from each pump pulse is annihilated and signal and idler photons, with carrier angular frequencies ω s and ω i , respectively, are created as a photon pair (conventionally ω s > ω i ). A temporal delay τ may be applied to pump 1 relative to pump 2. The carrier frequencies ω s and ω i are determined by the energy conservation constraint ω p1 + ω p2 = ω s + ω i as well as the phasematching conditions that are specific to the χ (3) medium. For this work we choose polarization maintaining fiber (PMF) as the nonlinear medium, where the polarization of both pump pulses are aligned with the slow axis of the PMF while the signal and idler photons are generated with polarization along the fast axis. In such a design the phasematching conditions are given by [18]: where k(ω) = n(ω)ω/c, n(ω) is the fast-axis (effective) refractive index in the fiber, which we model based on the Sellmeier equation of bulk silica [19], ∆n is the fiber birefringence, and c is the speed of light in vacuum. Starting with the |vac state, in which no photons exist in either signal or idler modes, the output state of the dual-pump SFWM process can be evaluated perturbatively as [16]: where |ν s , ν i represent a photon-pair state in which the signal (idler) angular frequency is ω s(i) + ν s(i) and κ is the interaction coupling constant that depends on the relevant χ (3) nonlinear susceptibility, the fiber length L, and pump powers, but not on the time delay τ that is applied to pump 1 relative to pump 2 before they are coupled into the fiber. The unnormalized joint spectral amplitude is given by [16,18]: Here σ 1(2) denotes the pump 1 (pump 2) spectral bandwidth (half width at 1/e 2 maximum amplitude); σ = σ 1 σ 2 / σ 2 ) is the group delay difference between the signal (idler) and the average group delay of the pumps acquired during the propagation in the fiber, k ′ s(i) = dk/dω| ω s(i) is the inverse group velocity of the signal (idler) in the fiber, τ p = L(k ′ p1 − k ′ p2 ) is the group delay between the two pumps acquired during the propagation in the fiber, and k ′ p1(2) = dk/dω| ω p1(2) +∆n/c is the inverse group velocity of pump 1 (2) in the fiber. The probability p(τ ) that a photon pair is generated, p(τ ) = |κ| 2 × dν s dν i |F (ν s , ν i ; τ )| 2 , is given by where p max corresponds to the maximum generation probability, which occurs when pump 1 and pump 2 maximally overlap in the middle of the fiber, i.e., τ = −τ p /2.
DEMONSTRATION AND CHARACTERIZATION OF PHOTON-PAIR GENERATION
We study the statistical properties of photon pairs generated in the dual-pump SFWM scheme using the experimental setup shown in Fig. 1(b). A Ti:sapphire modelocked laser with 80 MHz repetition rate, 772 nm central wavelength and 8 nm full-width-at-half-maximum (FWHM) bandwidth, pumps an optical parametric oscillator (OPO). The residual pump output from the OPO is used as pump 1. Pulses are generated in the OPO at central wavelengths between 530-660 nm with typical FWHM bandwidths of 1-3 nm, corresponding to pump 2. Pump 1 is time-delayed from pump 2 by an amount τ using an automated translation stage. The two pump paths are combined on a dichroic mirror and are coupled into a 1.6-cm long polarization-maintaining fiber (PMF) (PM630-HP, birefringence 3.5 × 10 −4 ), which serves as the nonlinear medium in which the SFWM process takes place [19][20][21]; the pump polarizations are aligned along the slow axis of the fiber. The signal and idler photons are produced with orthogonal polarization along the fast axis. A polarizer, which allows the signal and idler photons through, rejects most of pumps 1 and 2 and reduces noise from spurious interactions in the fiber. The signal and idler photons are separated by a dichroic mirror and each is coupled into a single-mode fiber connected to avalanche photodiode (APD) single-photon detectors (Excelitas SPCM-AQ4C). The detection signals, which are represented by electronic pulses from the APDs, are collected by a time-to-digital converter (TDC) that timestamps and processes detection events to record the number of single detection events C s , C i and C s ′ at APD s , APD i and APD s ′ , respectively; two-fold coincidences C si With 70 mW average power in pump 1 at 772 nm, and pump 2 set at 622 nm with 20 mW power, we record the number of detection events at the signal arm C s , idler arm C i , and coincidence counts C si , as a function of time delay τ between the two pump pulses. The results are presented in Figs. 2(a)-2(c). At a certain delay τ = τ 0 the counts reach peak values, as expected for dualpump SFWM, in which the photon-pair generation rate depends on the overlap between the two pump pulses [18]. For large |τ | (where the two pumps do not overlap and thus no dual-pump SFWM occurs) they asymptotically approach non-zero lowest values that amount to background photons and detection events, mainly due to Ra-man scattering, but also from single-pump SFWM, ambient light and dark counts. Figure 2(d) presents the cross-correlation g (2) si = C si R/C s C i , where R is the number of dual-pump pulse pairs over which the counts are taken (which is the laser repetition rate times the measurement duration). In order to ensure that the counts are synchronized with the laser pulses, the unconditional C s and C i counts in Figs. 2(a)-2(b) are gated at the laser repetition rate divided down to 8 MHz (due to bandwidth limitations of the TDC), and have been multiplied by 10 to reflect counts at the laser repetition rate; these values are used to calculate the cross-correlation g (2) si in Fig. 2(d). At the peak, g (2) si = 11.98 ± 0.02, indicating that correlated detection events of signal and idler photons occur. For large |τ |, g (2) si → 1, as expected from Poisson statistics of counts that originate from uncorrelated noise, confirming that signal-idler pairs indeed originate from dual-pump SFWM. We also measure the conditional auto-correlation of signal photons upon idler photon detection, g ss ′ |i = 0.017 ± 0.002, which indicates a low probability of multi-photon emission in one arm upon photon detection in the other arm.
Generally, the counts are given by: where N s (N i ) and η s (η i ) are the noise counts and the de-tection efficiency (accounting for both collection and de-tector efficiencies) of the signal (idler) photons generated in the dual-pump SFWM, respectively. In experiment, data is collected at various positions of the delay stage in Fig. 1(b). When the stage is in its central position, pump 1 is delayed relative to pump 2 by an unknown delay τ c ; the position of the stage is translated to create relative temporal delays τ exp of pump 1. Fitting curves to the data in Figs. 2(a)-2(c) are generated by substituting Eq. (4) into Eqs. (5) with τ = τ exp − τ c , and simultaneously fitting the three curves to the data, with N s , N i , η s , η i , p max , σ, τ p and τ c as common fitting parameters to all curves. The fitting result gives a maximal photon-pair generation probability per dual-pump pulse of p max = (6.0 ± 0.2) × 10 −3 , and collection efficiencies of η s = 13.4 ± 0.5% and η i = 10.7 ± 0.1%. The good agreement between model and experiment supports our approach, but we note that modeling is not required to determine the noise contributions, which can be obtained directly from a single measurement at large |τ | (where p(τ ) → 0), or three measurements of counts collected once when only pump 1 is present (no pump 2), once when only pump 2 is present (no pump 1) and once when both are blocked. Generally, quantifying noise enables one to gain information about source performance [5]; here, using Eqs. (5), knowledge of the noise enables us to extract the source performance from raw counts.
EFFECT OF PUMP DETUNING ON PHOTON-PAIR STATE AND SINGLE-PHOTON STATE PURITY
Dual-pump SFWM also provides enhanced capabilities in generating photon-pair quantum states with engineered properties [16,18]. Generally, the spectral quantum state of a photon pair can be expressed as |Φ = dν s dν i f (ν s , ν i ) |ν s , ν i , where f (ν s , ν i ) is the normalized joint spectral amplitude (JSA). The quantum state of the signal (idler) is then given by the density matrix ρ s(i) = Tr i(s) (|Φ Φ|), where Tr i(s) represents the partial trace over the idler (signal) degrees of freedom. The purity of the signal and idler photons P = Tr(ρ 2 s ) = Tr(ρ 2 i ) amounts to the degree to which they are in pure quantum states rather than mixed states, and is a critical figure of merit [22] in quantum protocols that rely on two-photon interference. Many efforts are being put into engineering the properties of photon pairs [23][24][25][26][27]. In particular, one of the most useful states is the factorable state, where the JSA can be written as independent wavefunctions of the signal (f s (ν s )) and idler (f i (ν i )) photons; that is, f (ν s , ν i ) = f s (ν s )f i (ν i ), leading to pure (P = 1) quantum states of signal and idler photons. Conversely, when the two photons are spectrally entangled (f (ν s , ν i ) is not factorable), P < 1 and the individual photons are in a mixed state.
Measurements of the joint spectral density
To characterize the JSA properties we measure the joint spectral density (JSD), |f (ν s , ν i )| 2 , using stimulated four-wave mixing, as proposed in [28] and demonstrated in [29,30] (experimental setup shown in Fig. 3). This is performed by adding a tunable Ti:sapphire continuouswave laser that co-propagates with the pumps in the PMF and seeds the idler beam to stimulate the creation of signal-idler pairs. Spectra are collected for each seed wavelength to generate the JSD. In the degenerate (single) pump case the JSA is given by [16]: where N is a normalization factor. The measured JSD for this case is shown in Fig. 4(a). It is evident that in addition to the main peak, there are sidelobes; these are due to the wings of the sinc function in f degen (ν s , ν i ), and originate from the sudden onset and ending of the nonlinear interaction when the pump pulse enters and exits the fiber. The fiber length L ≈ 1.6 cm in our experiments is chosen based on the model such that the purity of the photons is the highest it can be in the single-pump configuration, reaching a value of ∼ 83%. While this purity is high considering that no narrow spectral filtering is applied, the sidelobes seen in Fig. 4(a) -which constitute strong correlations between the signal and idler photons -limit the ability to achieve a factorable state [20].
In the dual-pump SFWM experiments, the time delay between pump 1 and 2 is set to be τ = τ 0 = −τ p /2 such that the two pumps maximally overlap at the center of the fiber and thus photon-pair production probability is highest. In this case Eq. (3) yields the JSA If the temporal walk-off between the pumps is large enough such that they completely sweep across each other within the medium, i.e. στ p ≫ 1, the . Under such conditions the JSA is expressed as the product of two Gaussian functions, which is the ideal expression for obtaining a factorable state [31] in general and possesses no sidelobes in particular. It becomes more factorable as the quantity C = σ 2 1 + σ 2 2 T s T i + (στ p ) 2 gets smaller, and, in principle, can become completely factorable when C = 0. We note that while the condition στ p ≫ 1 can be relatively easily satisfied by using a long medium, the value of C depends strongly on the dispersion characteristics of the medium. In PMF, στ p increases with detuning, while C decreases. Thence, we expect that increasing the detuning between the pumps would result in less correlations in the JSA [18].
With 20 mW average power each in pumps 1 and 2 and 30 mW average power in the seed beam, we obtain the experimental JSDs for various detunings ∆ = λ 1 − λ 2 (where λ 1(2) is the central wavelength of pump 1 (2)) shown in Fig. 4. As can be seen, with increased detuning the sidelobes' intensity weakens and the JSD becomes less correlated. Also shown in Fig. 4 (bottom row) are the corresponding calculated JSDs based on the model; the fidelity between measured (|f meas (ν s , ν i )| 2 ) and theoretical |f theory (ν s , ν i )| 2 JSDs [19] dν s dν i |f theory (ν s , ν i )| 2 |f meas (ν s , ν i )| 2 is > 95% for all measurements. These results support the models and the feasibility of the dual-pump approach to generating heralded single photons in pure wavepackets.
Purity measurements through autocorrelation
While the JSDs provide useful information about the photon-pair inter-correlations, they do not include details about the joint phase, and thus bear limited information about the individual photon purity. Photon-number statistics provide the purity [32] directly via the unconditional auto-correlation function [33] measured with the setup in Fig. 1(b), g (2) ss ′ | τ = C ss ′ (τ )R/C s (τ )C s ′ (τ ) = 1 + P meas (τ ), where P meas is the measured quantum-state purity. However, this kind of measurement is highly susceptible to noise contributions, which affect the detection statistics, resulting in an inaccurate deduced purity. Here again we find that the dual-pump scheme provides an advantage for quantifying the properties of the source. We derive upper and lower bounds for the purity of the signal photons P in the presence of noise. We assume two different types of noise: 1) Noise that is generated by the interaction of either pump and creates spurious photons at the signal arm together with an additional bosonthis boson could be another photon (e.g., through singlepump SFWM), or a collective excitation in the medium (such as a phonon). We call this type of noise spuri-ous noise. 2) Noise that occurs at the detector, such as dark counts or ambient light. We refer to this kind of noise as detection noise, with associated purity expressed through the auto-correlation function when both pumps are blocked, P det = (D ss ′ R/D s D s ′ − 1), where D s(s ′ ) is the detection-noise counts collected at APD s(s ′ ) and D ss ′ is the number of coincidences between APD s & APD s ′ , measured with blocked pumps. Given the raw purity P raw = P meas (τ 0 ) (since we are interested in measuring the purity when maximal photon-pair generation occurs), we can find bounds for the true purity of the signal photon produced through the dual-pump SFWM (see Appendix): where P noise = P meas (∞) is an effective purity associated with the total noise and is measured through the autocorrelation function at large |τ |; r = (1 − t s )(1 − t s ′ ) and t = √ t s t s ′ are the geometric averages of the ratios of signal-and noise-detections to the total counts, respectively, with t s(s ′ ) = C s(s ′ ) (∞)/C s(s ′ ) (τ 0 ); and u = √ u s u s ′ is the geometric average of the ratio of detection-noise to total noise, where u s(s ′ ) = D s(s ′ ) /C s(s ′ ) (∞). If P noise = 0 then P det = 0 necessarily; in such a case, or when t = 0 (no noise), the two bounds merge and the equality holds in Eqs. (8).
We measure P raw , P noise , r and t for various pump detunings. We find that P noise ∼ 0 for all measurements; we thus assume that it is zero and Eqs. (8) turn into the equality P = P raw /r 2 . The results of these measurements are summarized in Tab. I, together with a comparison to the theory [18]. The good agreement between P and the model, and the trend of improving purity with detuning, is yet another confirmation for the dual-pump approach as a superior technique to generate signal and idler pairs with each photon in a pure quantum state. The fact that we cannot use the above procedure to find r, t and P for the measurements at ∆ = 0 (single pump centered at 715 nm) emphasizes the advantage of the dual-pump scheme, where one can deduce the quantum state purity of the photons that are truly produced in pairs.
CONCLUSION
In conclusion, we experimentally investigate the generation of photon pairs through SFWM using two spectrally distinct laser pulses. We devise a new technique that utilizes the dual-pump nature for characterizing the performance of the photon-pair source in terms of generation probability, photon-collection efficiency, noise levels and quantum state purity of the individual photons.
As examples of potential applications of these capabilities, one can differentiate between degradation of the source medium, changes in the efficiency of collection and variations in ambient light; by scanning time delay, one can also characterize changes in pulse duration and modify the location of maximal pump overlap in the medium to avoid localized defects. Such tools may be especially useful in quantum applications where characterization of source performance and troubleshooting needs to take place periodically and remotely, especially in cases where the source needs to be placed in hardto-access locations such as space, or in a network with a vast number of sources. In addition, we show that large spectral detuning between the two pump pulses results in the generation of a highly factorable photon-pair state, with single-photon purities up to 97.4±1.7% as determined using dual-pump-enabled noise measurements, far exceeding those attainable with a single pump in the same generation medium.
To perform this first demonstration we choose PMF as the nonlinear medium due to its maturity as an SFWM photon-pair source [19-21, 30, 34, 35], the straightforwardness of the experimental setup and the simplicity of the model, which has a long track record of matching well with experimental results. It is expected, though, that more sophisticated media will be able to better exploit the dual-pump SFWM and overcome some of the issues found in PMF; for example, it has been proposed that with an adequately engineered birefringent medium, the two pumps could differ in polarization [36] rather than wavelength, thus avoiding the need for laser beams at two wavelengths. Also, the use of crystalline media where the Raman gain exhibit narrowband peaks (as opposed to silica) would enable the elimination of Raman background in the photon-pair spectrum and thus reduce noise levels significantly.
APPENDIX: DETERMINATION OF THE EFFECTIVE PURITY IN THE PRESENCE OF NOISE
In this Appendix we derive the inequalities in Eq. (8) that establish upper and lower bounds for the true quantum state purity of the signal and idler photons from measurements of the unconditional auto-correlation function g (2) ss ′ on the signal arm in the presence of noise.
Spurious noise
We first consider the effect of noise that is generated by the interaction of either pump and creates spurious photons at the signal arm together with an additional boson -this boson could be another photon (e.g., through single-pump SFWM), or a collective excitation in the medium (such as a phonon). We assume that before either of the pump pulses enters the fiber medium, the signal, idler, and any relevant collective excitation in the matter are in the vacuum state |vac . The final state after the two pump pulses leave the fiber is given by [32] where β and γ are the amplitudes of the dual-pump SFWM and noise generation, respectively,b † (Ω) is the creation operator of a boson with properties tagged by Ω and g(ν s , Ω) is the joint amplitude of the noise photon at the signal mode and the boson that is created in the interaction. We assume that all interactions are weak, i.e., |β| 2 ,|γ| 2 ≪ 1. The first order in γ, β is the lowest-order non-vacuum state, given by where are the states associated with the creation of a signal-idler pair through the dual-pump SFWM interaction (|ψ si ) and a signal-boson pair created by spurious processes (|ψ sb ). The signal density matrices associated with these states are ρ s = Tr i (|ψ si ψ si |) for the photon-pair state and ρ spu = Tr b (|ψ sb ψ sb |) for the spurious photons. The auto-correlation second-order coherence can be evaluated to yield [32,33] g (2) ss wherẽ P raw = w 2 P + (1 − w) 2 P spu + 2w(1 − w) Tr (ρ s ρ spu ) (13) is the raw measured purity, P = Trρ 2 s and P spu = Trρ 2 spu are the state purities of the signal and spurious photons, respectively, and w = |β| 2 /(|β| 2 + |γ| 2 ) is the ratio of the number of signal photons generated through dualpump SFWM to the total number of photons. Since 0 ≤ Tr(ρ s ρ spu ) ≤ P P spu [37], we can derive upper and lower bounds for the measured purity: When P spu = 0 or w = 0, 1, the two bounds merge and the equality holds.
Detection noise
Dark counts and ambient light that reaches the detectors constitute false detections that add background counts to the counts associated with photons that are created in the fiber. To model the effect of this type of noise we refer to the experimental setup in Fig. 1(b), where APD s and APD s ′ are used for signal auto-correlation measurements. Let us designate p s , p s ′ and p ss ′ as the probabilities of detection events at APD s , APD s ′ and coincidences between the two, respectively, after an interaction with a single pair of dual pumps, in the absence of detection noise. Similarly, we designate q s , q s ′ and q ss ′ as the probabilities of detecting noise events (which can be measured when both pumps are blocked) at APD s , APD s ′ or related coincidences between the two, respectively. All probabilities are assumed to be much smaller than 1, allowing perturbative calculations. By definition,g (2) ss ′ = p ss ′ /p s p s ′ = 1 +P raw . Similarly, we define P det = q ss ′ /q s q s ′ − 1. It follows then that the experimental auto-correlation, which includes photon-pair generation, spurious noise, and detection noise, is given by: g (2) ss ′ = p ss ′ + q ss ′ + p s q s ′ + p s ′ q s (p s + q s ) (p s ′ + q s ′ ) where v s(s ′ ) = q s(s ′ ) /(p s(s ′ ) + q s(s ′ ) ) is the ratio of detection noise counts to total counts on APD s(s ′ ) . Defining P raw = g (2) ss ′ − 1 as the raw purity measured in experiment that includes noise contributions, and using Eqs. (14), Eq. (15) turns into the inequalities When the two pump pulses are far delayed from each other, no dual-pump SFWM takes place, w = 0 and the equality holds for P raw ; we call the value of P raw in this case the "purity" of the total noise, designated as P noise = (1 − u s )(1 − u s ′ )P spu + u 2 P det , where u s(s ′ ) = v s(s ′ ) /((1 − w)(1−v s(s ′ ) )+v s(s ′ ) ) is the ratio of detection noise counts to the total noise counts on APD s(s ′ ) , and u = √ u s u s ′ . We further define r s(s ′ ) = w(1 − v s(s ′ ) ) as the fraction of dual-pump SFWM signal photon counts on APD s(s ′ ) to the total counts on APD s(s ′ ) , and t s(s ′ ) = 1 − r s(s ′ ) as the fraction of noise counts to the total counts on APD s(s ′ ) . The above inequalities then become: P raw ≥ r 2 P + t 2 (1 − u s ) (1 − u s ′ ) P spu + t 2 u 2 P det , P raw ≤ r 2 P + t 2 (1 − u s ) (1 − u s ′ ) P spu + t 2 u 2 P det where we have defined r = √ r s r s ′ and t = √ t s t s ′ . In terms of P noise the inequality can be rewritten as: P raw ≥ r 2 P + t 2 P noise , P raw ≤ r 2 P + t 2 P noise + 2rt P (P noise − u 2 P det ), (18) and therefore, yields the upper and lower bound of the true signal photon purity as (Eq. (8)) P ≤ P raw − t 2 P noise r 2 , P ≥ P raw − t 2 P noise r 2 − 2t r 2 P raw (P noise − u 2 P det ).
FUNDING INFORMATION | 7,142.8 | 2019-04-10T00:00:00.000 | [
"Physics",
"Engineering"
] |
High Yield Super-Hydrophobic Carbon Nanomaterials Using Cobalt/Iron Co-Catalyst Impregnated on Powder Activated Carbon
: Synthesis of super-hydrophobic carbonaceous materials is gaining a broader interest by the research community due to its versatile application in separation processes, special coating technologies, and membrane distillation. Carbon nanomaterials (CNMs) may exhibit stable superhydrophobic character due to their unique physio-chemical features which can be further controlled based on customer requirements by optimizing the process variables. This study deals with the application of a bimetallic catalyst composed of iron (Fe) and cobalt (Co) to synthesize CNMs from powder activated carbon as a precursor. The process parameters were optimized to ensure super-hydrophobic surfaces. Chemical vapor deposition was utilized for the growth of carbon nanomaterials. The impact of input variables on the desired output of yield and contact angle was analyzed. The chemical vapor deposition process was optimized using the response surface methodology based on Box-Behnken design. The proportion of the catalysts and reaction time were the three input explanatory variables whereas the desired response variables were selected as the carbon yield (CY) and contact angle (CA). The synthesized super-hydrophobic materials were characterized using field emission scanning electron microscopy (FESEM), transmission electron microscopy (TEM), Raman spectroscopy, thermogravimetric analysis (TGA), and contact angle analysis. The comprehensive statistical study of the results led to a significant model and optimization. The highest CY (351%) and CA (173 ◦ ) were obtained at the optimal loading of 2.5% Fe and 2% Mo with a reaction time of 60 min. The images obtained from FESEM and TEM revealed the presence of two types of CNMs including carbon nanofibers and multiwall carbon nanotubes. Thermogravimetric analysis was carried out to observe the temperature degradation profile of the synthesized sample. Raman spectroscopic analysis was also used to observe the proportion of ordered and disordered carbon content inside the synthesized samples. The improved catalytic super-hydrophobic carbon nanostructured materials production process proposed by this study assures the stability and high yield of the product.
Introduction
Production and applications of carbon-based nanomaterials (CNMs) occupies a very important place in nanotechnology research and development. CNMs with different forms including carbon nanotubes (CNTs), carbon nanofibers (CNFs), etc., with their unique physical and chemical properties are still believed to offer solutions for environmental and technical challenges [1][2][3][4].
Super-hydrophobic CNMs, which exhibit contact angles (CA) > 150 • , played serious roles in tackling many technical problems related to chemical and physical natures of surfaces and contacts. Therefore, super-hydrophobic CNMs have been implemented in various applications, including drug delivery materials [5], adsorbents [6], antifouling and self-healing membranes [7], and others [8,9]. The hydrophobicity of carbonaceous materials is mainly dependent on the roughness and surface chemistry of the synthesized samples [8]. CNMs' growth over other materials can influence the chemistry and roughness of the surface of the synthesized nanomaterial [9]. Thus, emphasis should be given to synthesize the carbon superstructures containing different types of hybrids to ensure superhydrophobic characteristics by optimizing its roughness and physiochemical properties.
Until recently, several methods have been used to synthesize CNMs including chemical vapor deposition method (CVD) [10], carbon arc discharge method (CA) [11], highpressure carbon monoxide conversion (HiPco) process [12], and pulsed laser vaporization technique (PLV) [13]. The superior quality of produced nanomaterials can be assured using CA and PLV methods. However, the application of CA method is restricted due to the high processing temperature of around 2700 • C needed to evaporate the carbon atoms from solid carbon sources; while the PLV method requires vacuum conditions and continuous graphite target replacement [14]. Thus, the scale-up of these production processes for commercialization purposes is difficult. Hence, for larger-scale production, the CVD method has gained much more attention than the others [15]. The CVD method is considered the most suitable synthesis method as it can ensure the product's high quality and quantity simultaneously [16]. Materials like silica [17,18], alumina [19,20], zeolite [21], and recently MgO [22,23] were used as supports for active metals to develop different types of nanostructured carbons like single walled nanotubes (SWNTs), multi walled nanotube (MWNTs), and nanofibers. Nevertheless, powder activated carbon (PAC) is considered to be the most suitable precursor in this regard due to its economic feasibility and unique features, including high thermal stability, high surface area, and prospect for chemical modifications by simple means [24,25]. Nano to micro dimensional carbon having different proportions of graphitic or disordered regions in carbon-carbon composites can induce superior properties based on the usage of suitable catalysts and process parameter optimization [26]. PAC has been used to synthesize carbon nanofibers (CNFs) using acetylene (C 2 H 2 ) and iron (Fe) catalysts in the CVD process [27,28]. On the other hand, modification of CNMs structure has been done by other researchers to achieve artificial super-hydrophobic surfaces. Optimization of the super-hydrophobic CNMs synthesis process could be hampered due to the aggregation nature of CNMs and usually, it comprises impurities [29]. The hybrid PAC-CNMs maintains the chemical compatibility between these two materials and combines the favorable characteristics of both.
Catalysts commonly used for CNTs growth are some transition metals like iron (Fe), cobalt (Co), and nickel (Ni) [27]. Lots of literature was already devoted to the CNT growth using different metals and their alloys [30]. Transition metals have empty 'd' shells which enable them to interact with hydrocarbons resulting in greater catalytic activity. The metallic particles used in the CVD process serve as seeding agents for the nanotubes, subsequently, they robustly control the configuration and quality of the finally developed materials. For the decomposition of hydrocarbons, transition metals require support for the successful growth of nanotubes. It was found that not only growth rate and diameter are highly influenced by the catalyst type and composition, but also the microstructure and morphology [31]. Moreover, the synthesis of CNMs from hydrocarbons is affected by the bimetallic catalysts' synergistic effects [32][33][34]. For example, the growth of CNMs was improved by using bimetallic catalysts such as Fe-Co, Co-Mo, and Fe-Ni [35][36][37][38][39]. All the above-mentioned catalysts (single or bimetallic catalysts) were successfully used with substrates other than PAC in most cases due to the difficulties in forming metallic nanoclusters on the surface of PAC with conventional methods.
Combining more than one catalyst at the same time affects the characteristics of the grown CNMs due to the interaction between the substrate surface and different metal clusters. In this study, the combination of bi-catalyst (Fe/Co), which deposited on PAC as a substrate, is considered to produce super-hydrophobic CNMs with high yield and effective performance. Moreover, the reaction is also investigated as an effective process parameter in thermal CVD to decompose acetylene as a carbon source at 650 • C. The optimum carbon yield (CY) and contact angle (CA) were the objective functions used for the production process while the catalyst's composition and reaction time were regarded as the explanatory variables in the response surface method (RSM) approach. The ultimate goal of this study is to produce super-hydrophobic CNMs to be utilized as material composites for several applications such as sorption [6], membrane distillation [40], separation of organic mixtures, purification of water by adsorption techniques, and catalysis [41].
Materials and Reagents
Iron nitrate Fe (NO 3 ) 3 ·9H 2 O, cobaltous nitrate hexahydrate Co(NO 3 ) 2 6H 2 O, PAC, and acetone were purchased from Sigma Aldrich, Kuala Lumpur, Malaysia. Acetylene gas (C 2 H 2 ), hydrogen gas (H 2 ), and nitrogen gas (N 2 ) were purchased from GasLink Industrial Gases SDN BHD, Selangor, Malaysia. Analytical grade reagents and chemicals were used in this research. Thus, no additional purification step was necessary for conducting the experiments.
Synthesis of Binary Metal Catalyst
The bimetallic catalysts (iron and cobalt) were first dissolved in acetone using their salts. Then the incipient wetness method was used to deposit the catalyst over the surface of the PAC. Based on Table 1, different weight percent (w/w%) of catalyst samples were prepared. The ratios between PAC and catalysts were calculated based on the experimental design matrix provided by the Box-Behnken design (BBD). According to Table 1, the catalyst solution was prepared by dissolving the desired amount of the catalysts in 5 mL of acetone solution. The mixture was kept inside the universal glass bottle and stirred until all the catalysts salts were completely dissolved.
Powdered Activated Carbon Impregnation
The metallic catalysts (iron and cobalt) were impregnated on PAC. After the catalysts were dissolved in acetone, PAC (2 g) was mixed with that catalyst-solvent mixture in a glass bottle. The glass bottles containing catalyst dissolved in acetone and PAC were sonicated at 40 kHz and 60 • C to evaporate the solvent. The mixture was further dried at 100 • C overnight using a conventional drying oven. After drying, the mixture was crushed to powder. The powder mixture thus obtained was stored using a desiccator to prevent moisture adsorption and sent for further characterization and application.
Synthesis of Carbon Nanostructured Materials (CNMs)
The catalysts of desired ratio was mixed with PAC and transferred to a ceramic boat (50 mm OD, 40 mm ID, 1500 mm L). This boat was then inserted inside the CVD reaction tube. The inert atmosphere was created by the flow of nitrogen (N 2 ) gas at a flow rate of 200 mL/min inside the CVD reactor. The sample was first heated at 350 • C for 2 h under this inert blanket condition to prevent burning and ash formation. Hydrogen gas (200 mL/min) was passed at 450 • C for 2 h inside the CVD reactor to ensure the reduction process of the calcined sample. The sample then was cooled down to room temperature. At this stage, the amount of moisture present in the sample was determined by weighing the sample (WC).
Under atmospheric pressure, the reduced catalyst mixture was placed inside the tubular ceramic reactor. Inside the reactor, a mixture of acetylene (50 mL/min) and hydrogen (200 mL/min) gas was flown. For facilitating the growth of CNMs, acetylene to hydrogen gas ratio was kept at 1:4 and the reaction was carried out at 650 • C for different reaction time as provided by the design matrix. After the reaction was completed, the reactor was cooled down using N 2 gas flow (200 mL/min). The weight (WP) of the synthesized CNMs was taken. The carbon yield was calculated using Equation (1).
where Wp and Wc are the weight of the sample after and before the reaction, respectively.
Equipment and Measurements
All the weighing measurements were recorded by using four digits weighing balance (HR-202i, Japan) with a measurement range between 0.001 and 220 g. PAC was mixed with the catalyst solution and the resultant mixture was placed inside the ultrasonic bath model (JAC 2010 P, Gyeonggi-Do, Korea) to ensure proper impregnation of metal catalysts onto the PAC substrate. The path is equipped with three levels of sonication, timer up to 99 min, and heater up to 90 • C. The drying of the PAC and CNMs samples was carried out in a drying oven (Model 600-Memmert, Büchenbach, Germany), where the maximum temperature can be set up to 220 • C. The CVD process was carried out in situ using an OTF-1200-80 mm dual-zone tube furnace for the CNMs growth. The tube furnace contains a fused quartz tube having the dimensions (OD: 80 mm; ID: 72 mm; length: 1000 mm). The heating area was covered with a resistance heating glass wool procured from Isolite Ceramic Fiber Sdn. Bhd., Johor, Malaysia.
Characterization of the synthesized samples was carried out to classify the type and shape of nanomaterials. Field emission scanning electron microscopic analysis (FESEM-Hitachi-SU8000, Ibaraki, Japan) was used to observe the morphology of the prepared samples. The aluminum stubs were coated by a platinum layer by sputtering and the synthesized sample was placed over it for FESEM analysis. Transmission electron microscopy (TEM) observations were made with a Hitachi-HT7700, Ibaraki, Japan microscope at 120 kV. Synthesized samples were mixed with acetone and the mixture was ultrasonicated. After ultrasonication, a drop of the sample was deposited over the copper grid supported perforated carbon film. The average diameter and particle size distribution of the nanomaterials were calculated using ImageJ software 1.8.0_112. Contact angle (CA) Processes 2021, 9, 134 5 of 19 measurement was carried out by KRUSS Goniometer (DSA100). A glass microscope slide (76 × 26 × 1.2 mm) covered with double-sided adhesive tape was used to measure the CA between the nanomaterials and water, wherein a water drop (4 µL) was placed onto the CNMs surface pasted on a tape. The average of triplicate measurements for each sample was taken. The oxidation behavior of the prepared nanomaterials was examined by thermal gravimetric analysis (STA-851, Mettler Toledo, Polaris Parkway Columbus, USA) at a temperature range of 25-800 • C, a heating rate of 10 • C/min, and an oxygen flow rate of 20 mL/min. The Raman spectra of the CNMs were obtained, where the laser power was kept constant at 100 using an Ar + laser (514 nm) and it was focused (50× objective) over the spot having a size of around 1.5-2.0 µm (Renishaw in Via, Sheffield, UK).
Response Surface Methodology and Process
The influence of three sovereign variables, Fe%, Co%, and reaction time on the two responses, carbon yield (CY) and contact angle (CA), were determined. Design-Expert V7.0 software was used and the Box-Behnken design (BBD) was employed by using the response surface methodology (RSM) approach. The influence of the main and combined variables on the desired responses chosen here was examined. A total of 17 experimental runs were proposed by DOE based on BBD with three center points. The independent variable ranges studied were Fe% (0-5%), Co% (0-4%), and reaction time (20-60 min), while gas ratios, temperature, and gas type were fixed at acetylene to hydrogen ratio of 1:4, 650 • C, acetylene gas, respectively. The complete design matrix conforming the actual experimental design and the responses are illustrated by Table 1. The accuracy of the developed models with data reproducibility was observed using the analysis of variance (ANOVA) test. The coefficient of determination (R 2 ) values was calculated to observe the competency level of the proposed model.
Results and Discussion
CA was selected mainly as an objective function response for the modeling and optimization process to investigate the effects of synthesis conditions (catalysts compositions and reaction time). This technique is unusual in such kinds of reactions where other researchers usually focus on the yield or the geometry of the nanoproducts, which has also been covered in the current work. In a previous study reported by our research group, an optimum Fe catalyst percentage weight of 5% was found to give a high CY when Fe was impregnated as a mono-catalyst on PAC [24].
Modeling and Statistical Analysis
After conducting the 17 experimental runs, the responses (CY and CA) (as shown in Table 1) were fitted to different statistical models. The models, including mean, linear, two-factor interaction (2FI), quadratic, and cubic polynomial models, are presented in Tables 2 and 3, respectively. It was evident that R 2 of the quadratic model, as well as its smallest F-value showed the closest correlation and the model was significant for CY and CA analyses as compared to that of the other models, as its p-value probability, Prob > F, was estimated to be less than 0.05. Thus, the quadratic model was chosen. Table 4 illustrates the ANOVA results obtained for CY% responses. As noticed from Table 4, the main effects on CY namely, Fe% (A), Co%(B), the reaction duration or time(C), the combined effect of Fe% and Co% (AB), the combined effect of Fe% and reaction time (AC), the combined effect of Co% and reaction time (BC), the quadratic effects of Fe% (A 2 ), the quadratic effects of Co% (B 2 ), and the quadratic effects of the reaction time (C 2 ), the interaction of the quadratic effects of the amount of Fe% and the amount of Co% (A 2 B), and the interaction of the quadratic effects of the amount of Co% and the amount of Fe% (B 2 A), were significant, as their Prob > F values were less than 0.05. Therefore, we can infer that the A, B, C, AB, AC, BC, A 2 , B 2 , C 2 , A 2 B, and B 2 A were the main determinants for CY. The polynomial equation for yield is given by Equation (2) The ANOVA results obtained for contact angles (CA) are illustrated in Table 5. It was observed that the major effects on CA, namely A, B, C, AB, AC, BC, A 2 , B 2 , C 2 , A 2 B, and B 2 A, were significant as their Prob > F values were less than 0.05. Thus, these effects were selected as the CA polynomial model parameters as given by Equation (3) Figure 1b shows a simulation model as an assessment pictorial representation between model values of CA predicted from the above Equation (3) and the experimental results, which demonstrates a good correspondence between the predicted and experimental values of CA.
Effect of Catalysts Composition on CNMs Yield and Contact Angle
The effect of catalysts loading on CNT yield and hydrophobicity was investigated. 3D-surface contour plots with surface mesh were plotted to observe the cumulative effects of the process parameters on CY and CA using Figure 2a,b, respectively. Figure 2 shows the effects of the catalyst
Effect of Catalysts Composition on CNMs Yield and Contact Angle
The effect of catalysts loading on CNT yield and hydrophobicity was investigated. 3D-surface contour plots with surface mesh were plotted to observe the cumulative effects of the process parameters on CY and CA using Figure 2a,b, respectively. Figure 2 shows the effects of the catalyst loading on CY and CA at constant reaction time (40 min). Five percent loading of Fe as a single catalyst was reported in previous work to produce an optimum yield of CNMs with a minimum average diameter [42]. Using the same Fe loading in the reported work gave a CY of 76.38% which is in agreement with the previously mentioned published work [42]. However, in the current study, it was found that combining 2% of Co with 2.5% of Fe catalysts gave higher CY of CNMs. This was observed in sample S4 which recorded the best CY of 351.3%. Catalyst loading higher than 2% of Co gave less CY, and this is attributed to the possible agglomeration of Co metal rather than the small amount of Co which increased the chance of the growth of CNMs, while further increase in Co led to decreasing in CY. Table 1 shows that the CY values do not follow any definite trend with the change of weight ratio of the catalyst. The two catalysts demonstrated distinct activities and tendencies with varying reaction conditions. The magnitude of CY depends on certain factors including particle size and capability of the catalyst to decompose C 2 H 2 , the solubility of carbon, rate of diffusion of carbon, the chemistry of the catalyst, and also on the consistent distribution of the metallic particles over the PAC surface [43,44]. Combination of Fe and Co as a bimetallic catalyst resulting in reducing the melting point of Fe-Co to much lower values compared to the particular melting points of Fe and Co alone [45]. In other words, this combination of metals could form an alloy with eutectic behavior. This is the reason for which the CY of Fe-Co bimetallic catalysts is much higher than the CY of the metal catalysts individually. The addition of Co to Fe catalyst plays an important role as it can increase the initial conversion of C 2 H 2 and restrain the fast catalyst deactivation. In addition, the formation of Fe-Co eutectic mixture, results in a reduction of melting point which affects the absorption of Fe on the surface of Co which consequently affects the shape of the catalyst, leading to the variation of the CNMs shape.
The hydrophobicity of the prepared CNM samples is expressed in terms of CA and presented in Table 1 and Figure 3. PAC is known to be a hydrophilic material, as PAC cast film shows a CA of 65 • , but the presence of CNMs which grew using 5% Fe as a catalyst shows an apparent CA of 176 • . By using Fe-Co as a bimetallic catalyst, the CA value started to increase following similar trends of CY, and for the best CY, the apparent CA is 173 • (Figure 3b). This suggests that the formation of more CNMs is the main reason for the high hydrophobicity. The main two factors affecting the surface wettability of CNMs are the surface roughness and surface chemistry [46,47]. The CNMs growth imparts a rougher surface area and reduces the gaps available for the droplets of water. Meanwhile, the higher CNMs growth leads to denser unfunctionalized stable carbon atoms which have a minimum affinity to attract water molecules. average diameter [42]. Using the same Fe loading in the reported work gave a CY of 76.38% which is in agreement with the previously mentioned published work [42]. However, in the current study, it was found that combining 2% of Co with 2.5% of Fe catalysts gave higher CY of CNMs. This was observed in sample S4 which recorded the best CY of 351.3%. Catalyst loading higher than 2% of Co gave less CY, and this is attributed to the possible agglomeration of Co metal rather than the small amount of Co which increased the chance of the growth of CNMs, while further increase in Co led to decreasing in CY. By using Fe-Co as a bimetallic catalyst, the CA value started to increase following similar trends of CY, and for the best CY, the apparent CA is 173° (Figure 3b). This suggests that the formation of more CNMs is the main reason for the high hydrophobicity. The main two factors affecting the surface wettability of CNMs are the surface roughness and surface chemistry [46,47]. The CNMs growth imparts a rougher surface area and reduces the gaps available for the droplets of water. Meanwhile, the higher CNMs growth leads to denser unfunctionalized stable carbon atoms which have a minimum affinity to attract water molecules.
Effect of Reaction Time on Yield and Contact Angle
The impact of the time duration for the CVD process on CY and CA was examined. The cumulative effects of the reaction time and catalyst composition on CY and CA were explicitly illustrated by 3D mesh diagrams. These diagrams were plotted based on the empirical model developed earlier, and are shown for CY (in Figure 4a,c) and for CA (in Figure 4b,d). Figure 4a,c show the response surface plots for the effects of the reaction time at a fixed 2% of Co on: (a) CY and (c) CA, and at fixed 2.5% of Fe% on: (b) CY and (d) CA. There was a significant increment in CY with increasing reaction time, from 20 to 40 min and then to 60 min, as shown in Figure 4a,c. The reaction time is the dominating parameter for CY and the highest CY was obtained at a reaction time of 60 min. The increased CY with increasing the reaction time can be considered as an indication for the catalyst stability which remains constantly active in performing mixed with that catalyst-solvent mixture C2H2 decomposition. Eventually, the constant activity of the catalysts and high carbon moieties diffusion on the PAC resulted in more accumulative CNMs production. This reveals the general understanding of the interaction between the most effective parameters in the process,
Effect of Reaction Time on Yield and Contact Angle
The impact of the time duration for the CVD process on CY and CA was examined. The cumulative effects of the reaction time and catalyst composition on CY and CA were explicitly illustrated by 3D mesh diagrams. These diagrams were plotted based on the empirical model developed earlier, and are shown for CY (in Figure 4a,c) and for CA (in Figure 4b,d). There was a significant increment in CY with increasing reaction time, from 20 to 40 min and then to 60 min, as shown in Figure 4a,c. The reaction time is the dominating parameter for CY and the highest CY was obtained at a reaction time of 60 min. The increased CY with increasing the reaction time can be considered as an indication for the catalyst stability which remains constantly active in performing mixed with that catalyst-solvent mixture C 2 H 2 decomposition. Eventually, the constant activity of the catalysts and high carbon moieties diffusion on the PAC resulted in more accumulative CNMs production. This reveals the general understanding of the interaction between the most effective parameters in the process, namely catalysts composition and reaction time. The constant catalysts activity which resulted in constant increasing of CNMs production can open a window for contentious CNMs mass production.
The effects of the reaction time on CA at a fixed amount Fe% (2.5%) and Co% (2%) are shown in Figures 4b and 4d, respectively. There was a noteworthy enhancement of CA with enhancing reaction time, starting from 20 to 60 min, as reflected by both plots. Although CA increases with increasing reaction time, CA has a limitation of 180 • maximum. As such it is not expected to have a CA beyond the maximum limit even though the reaction time is increased.
namely catalysts composition and reaction time. The constant catalysts activity which resulted in constant increasing of CNMs production can open a window for contentious CNMs mass production.
The effects of the reaction time on CA at a fixed amount Fe% (2.5%) and Co% (2%) are shown in Figure 4b,d, respectively. There was a noteworthy enhancement of CA with enhancing reaction time, starting from 20 to 60 min, as reflected by both plots. Although CA increases with increasing reaction time, CA has a limitation of 180° maximum. As such it is not expected to have a CA beyond the maximum limit even though the reaction time is increased.
Optimization of Chemical Vapor Deposition (CVD) Process
After obtaining the significant empirical statistical model, it is possible to optimize the conditions which enable designing the desired combination of all factors. From the DoE software, a numerical optimization algorithm was selected and used to optimize the unconstrained variables and desired responses. All the factors and responses corresponding to the upper and lower limits of the experimental range had to satisfy the optimization criteria which defines the desired limitations and constraints. In this study, C2H2 decomposition into CNMs was maximized to obtain the CY and CA values with the importance of five out of five, while the other parameters were kept in the chosen experimental range. Some sets of predicted solutions were obtained as potential optimization conditions, which were additionally classified by desirability as listed in Table 6. The highest
Optimization of Chemical Vapor Deposition (CVD) Process
After obtaining the significant empirical statistical model, it is possible to optimize the conditions which enable designing the desired combination of all factors. From the DoE software, a numerical optimization algorithm was selected and used to optimize the unconstrained variables and desired responses. All the factors and responses corresponding to the upper and lower limits of the experimental range had to satisfy the optimization criteria which defines the desired limitations and constraints. In this study, C 2 H 2 decomposition into CNMs was maximized to obtain the CY and CA values with the importance of five out of five, while the other parameters were kept in the chosen experimental range. Some sets of predicted solutions were obtained as potential optimization conditions, which were additionally classified by desirability as listed in Table 6. The highest desirability for the optimum process conditions solution was 2.73% of Fe, 2.42% of Co, and 60 min reaction time, which predicted to give CY of 351.30% and CA of 175.88 • . For experimental verification, an experiment was conducted at the suggested optimized conditions. These results show that the experimental values were CY of 349.54% and CA of 173 • , which are consistent with the model predicted values.
Surface Morphology Analysis
The morphological features of the synthesized samples were analyzed using FESEM technique. Figure 5a-f showed the images of the prepared CNMs for six selected samples (S1, S4-S6, S10, and S11, respectively) with 50,000× magnification. Figure 5c displayed the FESM image of S5 which was found to have the highest CY and CA by using 5% of Fe as a mono-catalyst and reaction time of 40 min. Two types of carbon nanomaterials were found in this sample as carbon nanofibers (CNFs) and multiwall carbon nanotubes (MWCNTs). Figure 5f displayed the FESEM image of S11 which was prepared by using 4% of Co as a mono-catalyst and reaction time of 40 min. It was found that using a high amount of Co only as a catalyst gave very poor CY which indicates very poor growth of CNMs. Figure 5e displays the FESEM image of S10 which was prepared by using 5% of Fe and 4% of Co as bimetallic catalysts and reaction time 40 min. This image revealed that using a high amount of Co and Fe produced mainly aggregated CNMs and CNFs, as such experiment conditions are not preferred. Figure 5a,b, and d displayed the FESM image of S1, S4, and S6, respectively. These images show three different types of CNMs as helix-like CNFs, CNFs, and MWCNTs. It was noticed that increment of reaction time did not alter the shape of the CNMs and it affected only the CY of the CNMs.
Surface Morphology Analysis
The morphological features of the synthesized samples were analyzed using FESEM technique. Figure 5a-f showed the images of the prepared CNMs for six selected samples (S1, S4-S6, S10, and S11, respectively) with 50,000× magnification. Figure 5c displayed the FESM image of S5 which was found to have the highest CY and CA by using 5% of Fe as a mono-catalyst and reaction time of 40 min. Two types of carbon nanomaterials were found in this sample as carbon nanofibers (CNFs) and multiwall carbon nanotubes (MWCNTs). Figure 5f displayed the FESEM image of S11 which was prepared by using 4% of Co as a mono-catalyst and reaction time of 40 min. It was found that using a high amount of Co only as a catalyst gave very poor CY which indicates very poor growth of CNMs. Figure 5e displays the FESEM image of S10 which was prepared by using 5% of Fe and 4% of Co as bimetallic catalysts and reaction time 40 min. This image revealed that using a high amount of Co and Fe produced mainly aggregated CNMs and CNFs, as such experiment conditions are not preferred. Figure 5a,b, and d displayed the FESM image of S1, S4, and S6, respectively. These images show three different types of CNMs as helix-like CNFs, CNFs, and MWCNTs. It was noticed that increment of reaction time did not alter the shape of the CNMs and it affected only the CY of the CNMs. S5, (d) S6, (e) S10, (f) S11. Figure 6 shows the TEM images of S4. The images revealed the MWCNTs internal structure and helix-like CNFs as well as CNFs, which were attached over the surface of amorphous PAC. The internal cavity or hollow places represent the development of MWCNTs, while the other CNMs were observed as a solid structure. Some darker spots observed inside the TEM images represent the bimetallic nanocluster which initiates the buildup of carbon atoms resulting in CNTs and other CNMs formation. It is also clear from the position of the catalyst nanocluster on the top end of the CNMs that the growth mechanism could be following the top-down model. This is in agreement with our findings that CY showed an increasing trend with increasing reaction time. In that case, the bimetallic catalyst nanoclusters were not covered. Thus, their exposure enables them to be catalytically active and decompose C2H2 into CNMs. In addition, TEM images exhibited a crooked or twisted CNMs layout. According to the proposed catalytic growth mechanism, the crooked, or twisted CNMs may be the result of a variety of clusters of catalyst crystallization and morphology when conjugated with carbon segregation on the active sites around the catalyst periphery during CNMs growth. Figure 6 shows the TEM images of S4. The images revealed the MWCNTs internal structure and helix-like CNFs as well as CNFs, which were attached over the surface of amorphous PAC. The internal cavity or hollow places represent the development of MWCNTs, while the other CNMs were observed as a solid structure. Some darker spots observed inside the TEM images represent the bimetallic nanocluster which initiates the buildup of carbon atoms resulting in CNTs and other CNMs formation. It is also clear from the position of the catalyst nanocluster on the top end of the CNMs that the growth mechanism could be following the top-down model. This is in agreement with our findings that CY showed an increasing trend with increasing reaction time. In that case, the bimetallic catalyst nanoclusters were not covered. Thus, their exposure enables them to be catalytically active and decompose C 2 H 2 into CNMs. In addition, TEM images exhibited a crooked or twisted CNMs layout. According to the proposed catalytic growth mechanism, the crooked, or twisted CNMs may be the result of a variety of clusters of catalyst crystallization and morphology when conjugated with carbon segregation on the active sites around the catalyst periphery during CNMs growth. The ImageJ ® image processing software was used to determine the average diameter and particle size distribution of CNMs to perform statistically reliable measurements; several micrographs have been processed for each sample. The minimum number of investigated particles was 200. It was found that the average size diameter (as shown in Table 7) was the smallest when using Fe as a monocatalyst as in S5. Using bimetallic catalysts of Fe and Co tends to increase the average size diameter as Co catalyst percentage weight and reaction time increased. The ImageJ ® image processing software was used to determine the average diameter and particle size distribution of CNMs to perform statistically reliable measurements; several micrographs have been processed for each sample. The minimum number of investigated particles was 200. It was found that the average size diameter (as shown in Table 7) was the smallest when using Fe as a mono-catalyst as in S5. Using bimetallic catalysts of Fe and Co tends to increase the average size diameter as Co catalyst percentage weight and reaction time increased. The particle size distribution test result is depicted in Table 7. It was found that using Fe as a mono-catalyst gave around 64% and 20% in the range between 0-50 and 50-100 nm, respectively. However, adding a small amount percentage weight of Co increased the size of CNMs. For example, S4 gave around 21% and 47% in the range between 0-50 and 50-100 nm, respectively. It is worth mentioning that the best samples satisfying the definition of CNMs (<100 nm) are those with 2% Co and 2.5% Fe.
Raman Analysis
Raman spectroscopy was used to characterize the nature of the prepared CNMs as shown in Figure 7 for S4 and S5. There was an observation for two well-unattached peaks at~1590 cm −1 for G peaks which corresponds to the movement of the two adjacent carbon atoms moving in opposite directions inside the graphite sheet [48], and at~1350 cm −1 for D peaks which corresponds to sp3-hybridization of carbon atoms at the sidewalls of CNTs [49]. These peaks usually appear in multiwall CNTs, while the radial breathing mode (RBM) can be spotted in the range of 100-400 cm −1 in SWCNTs. No (RBM) peaks were observed for the samples which indicates that no SWCNT was produced. The ID/IG ratios were calculated to estimate the variation of CNMs crystallinity with Fe and Fe-Co catalysts. The ID/IG for the S5 sample is 0.63 which was lower than those using Fe-Co bimetallic catalysts for the S4 sample (0.96). This suggests that the use of the Fe-Co bi-catalysts system will provide the sample containing well developed graphitic structure than that of the sample using the Fe catalyst only [50].
SWCNTs. No (RBM) peaks were observed for the samples which indicates that no SWCNT was produced. The ID/IG ratios were calculated to estimate the variation of CNMs crystallinity with Fe and Fe-Co catalysts. The ID/IG for the S5 sample is 0.63 which was lower than those using Fe-Co bimetallic catalysts for the S4 sample (0.96). This suggests that the use of the Fe-Co bi-catalysts system will provide the sample containing well developed graphitic structure than that of the sample using the Fe catalyst only [50].
Thermal Stability Analyses
Thermal stability, quality, and purity of the synthesized CNMs were analyzed using TGA analysis. The curve obtained for TGA and derivative thermogravimetric (DTG) analysis of samples S4 and S5 are illustrated by Figure 8. TGA analysis demonstrates the weight loss at different temperatures which reflects the thermal stability of the finally obtained CNMs. The first observation
Thermal Stability Analyses
Thermal stability, quality, and purity of the synthesized CNMs were analyzed using TGA analysis. The curve obtained for TGA and derivative thermogravimetric (DTG) analysis of samples S4 and S5 are illustrated by Figure 8. TGA analysis demonstrates the weight loss at different temperatures which reflects the thermal stability of the finally obtained CNMs. The first observation of weight loss was at a temperature of around 100 • C for both samples which can be attributed to the adsorbed moisture and hydrocarbons over the surface of synthesized CNMs [51]. Oxidation began at approximately 400 • C and 500 • C for sample S5 and S4, respectively which resulted in a loss of nearly 85 wt%. The oxidation temperature of different carbon structures is not the same and varies according to the inter-atom bonds' strength and their network uniformity. The TGA profile of S4 showed more thermal stability due to the presence of more CNMs. It was also noticed that no further weight loss after a temperature of 600 and 667 • C for sample S5 and S4, respectively. The residual materials obtained for the sample contain ash residues and it was estimated to be 6.5% and 8% for sample S5 and S4, respectively. This might be due to the catalytic impregnation of PAC earlier in the CVD process for CNMs growth. of weight loss was at a temperature of around 100 °C for both samples which can be attributed to the adsorbed moisture and hydrocarbons over the surface of synthesized CNMs [51]. Oxidation began at approximately 400 °C and 500 °C for sample S5 and S4, respectively which resulted in a loss of nearly 85 wt%. The oxidation temperature of different carbon structures is not the same and varies according to the inter-atom bonds' strength and their network uniformity. The TGA profile of S4 showed more thermal stability due to the presence of more CNMs. It was also noticed that no further weight loss after a temperature of 600 and 667 °C for sample S5 and S4, respectively. The residual materials obtained for the sample contain ash residues and it was estimated to be 6.5% and 8% for sample S5 and S4, respectively. This might be due to the catalytic impregnation of PAC earlier in the CVD process for CNMs growth.
Conclusions
This research deals with the dense growth of CNMs from abundantly available activated carbon using a binary metal mixture in a certain proportion using the CVD process. The RSM statistical
Conclusions
This research deals with the dense growth of CNMs from abundantly available activated carbon using a binary metal mixture in a certain proportion using the CVD process. The RSM statistical method was successfully used to show the interactions of different input variables for CNMs growth on PAC. The output responses including CY and CA were determined and significant regression models were obtained. Eventually, this led to the successful optimization of the CVD process. The proportion of Fe and Co catalysts, as well as reaction time, plays a crucial role in obtaining the highest CY and CA values. Process parameter optimization revealed that the mixture of 2.5% Fe and 2% Co is the best catalyst for CNMs growth when the CVD process is carried out for 60 min to have the highest magnitude of CY around 351% with CA values of 173 • . Surface morphological features observed by using SEM and TEM analyses demonstrated that mixtures of CNMs such as helix-like CNFs, CNFs, and CNTs were produced under optimum conditions. High degrees of graphitization with some defects inside the nanocarbon matrix was detected during Raman analysis. Thus, a stable, super-hydrophobic nanostructured carbon having different shapes had been successfully synthesized in this study from powdered activated carbon. The catalytic approach used here ensured high yield, CY% making the overall process economically feasible.
Conflicts of Interest:
The authors declare no conflict of interest. | 9,548 | 2021-01-11T00:00:00.000 | [
"Chemistry"
] |
Dynamic Sensing Performance of a Point-Wise Fiber Bragg Grating Displacement Measurement System Integrated in an Active Structural Control System
In this work, a fiber Bragg grating (FBG) sensing system which can measure the transient response of out-of-plane point-wise displacement responses is set up on a smart cantilever beam and the feasibility of its use as a feedback sensor in an active structural control system is studied experimentally. An FBG filter is employed in the proposed fiber sensing system to dynamically demodulate the responses obtained by the FBG displacement sensor with high sensitivity. For comparison, a laser Doppler vibrometer (LDV) is utilized simultaneously to verify displacement detection ability of the FBG sensing system. An optical full-field measurement technique called amplitude-fluctuation electronic speckle pattern interferometry (AF-ESPI) is used to provide full-field vibration mode shapes and resonant frequencies. To verify the dynamic demodulation performance of the FBG filter, a traditional FBG strain sensor calibrated with a strain gauge is first employed to measure the dynamic strain of impact-induced vibrations. Then, system identification of the smart cantilever beam is performed by FBG strain and displacement sensors. Finally, by employing a velocity feedback control algorithm, the feasibility of integrating the proposed FBG displacement sensing system in a collocated feedback system is investigated and excellent dynamic feedback performance is demonstrated. In conclusion, our experiments show that the FBG sensor is capable of performing dynamic displacement feedback and/or strain measurements with high sensitivity and resolution.
Introduction
Fiber Bragg grating (FBG) sensors possess many excellent properties such as small size, mass-production at low cost, and immunity to electro-magnetic interference (EMI). For more than a decade, they have been proven to be sensitive to many physical quantities such as strain, temperature, pressure, acceleration, and force [1][2][3]. Since many FBGs can be inscribed into a single fiber, they also have multiplexing ability to detect several different positions in the same structure simultaneously. Traditionally, FBG sensors are mounted on or imbedded in structures. Thus, detecting out-of-plane point-wise displacement can't be achieved without work-around methods such as bonding an FBG sensor to a cantilever structure to indirectly measure the displacement responses [4,5]. However, indirect sensing methods are still not capable of measuring point-wise displacement. Due to the different modes of the cantilever structures, indirect sensing methods are also not suitable for dynamic measurements.
To directly utilize an FBG to measure out-of-plane point-wise displacement, a method that can point-wisely glue an FBG sensor perpendicular to the surface containing the detecting point is adopted in this paper [6]. Since sensors are one of the key elements in smart structures [7], the purpose of this work is to investigate the feasibility of integrating the proposed FBG displacement sensing system into smart structures for performing active vibration suppression. A smart cantilever beam actuated by a piezoelectric actuator is employed to demonstrate the dynamic sensing ability as well as the feedback sensing ability of the proposed FBG displacement sensing system.
The sensing principle of the FBG is based on the shift of the Bragg wavelength resulting from variations of environmental physical quantities. Under static or low frequency condition, an optical spectrum analyzer (OSA) is usually used to detect wavelength shift [8,9]. However, since the demodulating speed of the OSA is limited, wavelength shifts are not able to be recorded in time for high frequency signals. On the contrary, demodulating techniques which transfer wavelength shifts to energy variations through optical filters are capable of detecting high frequency responses. Optical filters such as the long-period fiber grating (LPFG) filters, FBG filters, and chirped FBG filters are capable of providing different dynamic sensing ranges and resolutions [10][11][12]. Among above mentioned filters, the FBG filter can provide the highest sensitivity and it possesses high signal-to-noise (SNR) ratio due to its smallest full-width at half maximum (FWHM) compared to other grating-based filters. Thus, different from the sensing system in [6] which employs an LPFG filter to demodulate the displacement responses, this study employs an FBG filter in the sensing system to enhance the dynamic sensing performance.
There are several control algorithms that have been successfully utilized to control flexible smart structures [13]. Since the purpose of this study is focused on the dynamic sensing performance of the proposed point-wise FBG displacement sensor system in an active structural control application, a velocity feedback control is adopted in this paper to add damping of the smart cantilever beam and suppress the vibrations. The proposed FBG displacement sensor and the piezoelectric actuator is set up collocated to each other on top and bottom surfaces of the cantilever beam, respectively. Before performing experiments, an optical full-field measurement technique called amplitude-fluctuation electronic speckle pattern interferometry (AF-ESPI) is used to provide full-field vibration mode shapes and resonant frequencies. To verify the dynamic demodulation performance of the FBG filter, a traditional FBG strain sensor calibrated with a strain gauge is first employed to measure the dynamic strain of impact-induced vibrations. Then, system identification of the smart cantilever beam is performed by FBG strain and displacement sensors. To verify dynamic sensing ability of the proposed FBG displacement sensor, a laser Doppler vibrometer (LDV) is simultaneously employed as a comparison for the displacement measurement. Finally, since the effect of the velocity feedback for sinusoidal waveforms is equivalent to delay the waveforms by some phases, a delay controller is also utilized in our work as a comparison for the velocity feedback controller. To our knowledge, this is the first time that the proposed FBG displacement sensing system being integrated into a smart structure for performing active vibration control.
This paper is organized as follows. Section 2 summarizes the sensing principle, calibration method, and set-up method of the FBG displacement sensor. The model of the smart cantilever beam is presented in Section 3. The velocity feedback controller and the delay controller are briefly introduced in Section 4. Finally, the experimental setup, performances of the FBG-filter based demodulation technique, system identifications performed by FBG strain/displacement sensors, and control results of suppressing vibrations of cantilever beam are reported and discussed in Section 5.
FBG Sensing System
A fiber Bragg grating (FBG) is a periodic distribution of the refractive index along the fiber core. From the Bragg's law, the Bragg wavelength λ S of an FBG sensor is given by [14]: where Λ is the Bragg grating period and eff n is the effective refractive index of the fiber core. The shift in Bragg wavelength due to an applied strain can be expressed as: where ε is the strain induced in the fiber, p e is the effective photoelastic coefficient, and λ S0 is the Bragg wavelength of the grating without the strain field. In order to measure the point-wise displacement of a point on a structure subjected to dynamic loadings using an FBG, one end of a fiber (of length l 0 ) containing the FBG needs to be fixed to a stationary boundary while the other end is fixed to the sensing point. The dynamic displacement of the sensing point on the structure can be obtained from the following relation: The spectrum of the light reflected by the FBG sensor can be approximated by a Gaussian function, given by: where R S is the maximum reflectivity of the FBG sensor, σ S is the grating full-width at half maximum (FWHM). In order to enhance the signal-to-noise (SNR) ratio and the sensitivity, an FBG filter located at the output of a broadband source (BBS) is used as a demodulator in the sensing system. If the intensity of the BBS is P(λ), the total light power I detected by the photodetector (PD) can be expressed as: where F(λ) is the spectrum of the FBG demodulator. For a broadband source with a relative flat spectrum, the power spectrum may be assumed to be a constant as P λ over the operation range. The PD transforms the light intensity to voltage signals. Before performing the measurement, calibration of the FBG sensing system is necessary to ensure the linearity, maximum sensing output, and sensitivity. The calibrating criteria are based on the demodulation behavior of the FBG filter. The transmittance of the FBG filter can be approximated by [14][15][16]: where λ F is the Bragg wavelength of the filter, R F is the maximum reflectivity, and σ F is the grating FWHM of the FBG filter. Substituting Equations (4) and (6) into Equation (5), the light intensity after evaluating the integration is expressed as: In Equation (7) and δλ nor is a normalized wavelength mismatch defined as: where S F S 0 F S δλ λ λ λ λ λ = − = − +Δ is the wavelength mismatch. Hence, the variation of intensity I will depend on the wavelength mismatch. It is noted that if λ F of the FBG filter attached to a translation stage is adjusted to the left of λ S before the measurement, the PD output electrical signal will increase as the FBG sensor is elongated and decrease as the FBG sensor is compressed. For the case that small dynamic strain is applied, the response of the FBG sensor is linear and the wavelength-to-intensity conversion factor K depending on the operation point can be derived from Equation (7) The maximum wavelength-to-intensity conversion factor K nor,max can be obtained by setting dK nor /d(δλ nor ) = 0, where K nor = K/2I P . The normalized wavelength mismatch at K nor,max is 1 / 2 which indicates that an optimal operation point of δλ is: If the optimal operation condition, i.e., δλ OPT , is matched, then the maximized PD output can be obtained. Meeting of the optimal operation condition is a method to calibrate the FBG sensor before the experiment. An optical spectrum analyzer (OSA) can be employed to monitor the optimal operation condition during the experiment. In this paper, the proposed FBG displacement sensor with high sensitivity is used to detect and feedback the displacement responses of a smart cantilever beam. Since the cantilever beam is actuated by a piezoelectric actuator near the fixed end of the cantilever beam, we can apply a sinusoidal voltage input to the piezoelectric actuator at the first resonant frequency of the cantilever beam to excite the vibration and slightly adjust the translation stage of the FBG filter to prevent the displacement waveform from saturation.
To set up an FBG as an out-of-plane FBG displacement sensor, one end of the FBG displacement sensor is glued to a vertical translation stage and the other end of the FBG is glued to the sensing point on the surface of the cantilever beam with a mix of epoxy resin and hardener. Since FBGs are very sensitive to the longitudinal strain, the set-up method mentioned above allows an FBG to be capable of measuring the out-of-plane dynamic displacement of the cantilever beam with high sensitivity. The illustration of the set-up method for the proposed FBG displacement sensor is shown in Figure 1. Details regarding the FBG displacement sensor set-up method can be found in Chuang and Ma [6].
Model of the Smart Cantilever Beam
In this section, the model of the smart cantilever beam is briefly presented [17][18][19]. Due to the fact that the proposed point-wise FBG displacement sensor is very sensitive, it is attached near the fixed end of the cantilever beam to avoid saturations of measurement results. Since the thickness is small compared to the length and width of the cantilever beam, the effects of shear deformation and rotary inertia can be neglected. First, consider the free strain ε a of the actuating layer when a voltage v a is applied to a piezoelectric material polarized in z direction as: where d 31 is the piezoelectric constant of the piezoelectric material, t a is the thickness of the piezoelectric actuator. Assume bonding between the actuator and the beam is perfect and the strain distribution is linear across the thickness of the beam, the strain distribution can be decomposed into flexural component and longitudinal component as: The stress distribution inside the piezoelectric actuator is: where E a is the Young's modulus of the piezoelectric material. The stress distribution inside the beam is: where E b is the Young's modulus of the beam. By applying the moment equilibrium about the center of the beam: a t t y y y y y y σ σ + − + = (17) and the force equilibrium along the longitudinal axis of the beam: we can obtain ε 0 and a, in which: where t b is the thickness of the beam. The induced moment ( , ) M x t in the beam can be expressed as: where u stands for the unit step function, 1 x and 2 x are coordinates of the left and right ends of the piezoelectric actuator, and I b is the moment of inertia of the beam and: From the Euler-Bernoulli beam theory, the governing equation of motion of the cantilever beam actuated by a piezoelectric patch is given by: where w(x,t) is the out-of-plane displacement of the beam from its equilibrium, ρ is the mass density, and A b is the cross-sectional area of the beam. The boundary conditions of the cantilever beam are given by: where L is the length of the cantilever beam. The out-of-plane displacement w(x,t) can be expanded as an infinite series of eigenfunctions in the form: is the mode shape, also known as eigenfunction, satisfying the ordinary differential equations. From the following orthogonality properties: where ij δ is the Kronecker delta function. From the boundary conditions of the cantilever, the mode shape is given by: where i β are the roots of the following characteristic equation: The resonant frequencies i ω of the smart cantilever beam are evaluated from Multiplying the governing equation of the Euler-Bernoulli beam by ( ) j x φ and integrating over the length of the beam, we have: From the orthogonality properties Equations (25) and (26), the property of Dirac delta function, and with the assumption that the voltage v a (x,t) is constant in the range 1 2 [ , ] x x , we have: By adding natural damping of the beam, the governing equation ( ) i q t for the smart cantilever beam can be written as follows: In this study, a simple but effective negative velocity feedback controller is employed to suppress vibrations of the smart cantilever beam as well as to demonstrate dynamic sensing performance of the proposed FBG displacement sensor system. In velocity feedback control, the displacement responses obtained by the FBG sensor are differentiated and fed back to the piezoceramic actuator. The structure of the velocity feedback controller can be written as follows: where G is the constant control gain. Equation (31) can then be rewritten as: Thus, it is obvious that velocity feedback control has the effect to add damping of the flexible structure and can be used to suppress the vibration.
In this study, the velocity feedback controller is developed in Simulink and implemented by a dSPACE DS1104 system with a sampling frequency set to 50 kHz. Figure 2 shows the Simulink program in which the exciting and control process are combined into one program. First the smart cantilever beam is excited by the piezoceramic actuator to the steady state at the first or second natural frequency. During the steady-state vibration, actuator's input is terminated suddenly to induce free vibration. At the same time, feedback loop is connected to apply control algorithms to suppress the vibration of the cantilever beam. The cantilever beam used in the experiment is made of 6,061 aluminum. The dimensions of the cantilever beam are shown in Figure 1. The actuator used in this study is an APC 855 piezoceramic plate and its parameters are shown in Table 1. Before performing the sensing and feedback experiments, we employ an optical amplitude-fluctuation electronic speckle pattern interferometry (AF-ESPI) [20] technique to provide full-field vibration mode shapes as well as the resonant frequencies of the smart cantilever beam. Compared to the ordinary time-averaging method, the fringe patterns obtained by the AF-ESPI method largely enhance visibility and reduce noise. When the vibration frequency of the smart cantilever beam is near the resonant frequency, stationary distinct fringe patterns will be observed in a monitor. Thus, the AF-ESPI can be used to simultaneously obtain full-field vibration mode shapes as well as resonant frequencies. The mode shapes and resonant frequencies of the smart cantilever beam obtained by the AF-ESPI technique are shown in Tables 2 and 3 for the bending modes and the torsional modes, respectively. The brightest lines in AF-ESPI measurements are the nodal lines of the vibration mode shapes. In this study, a non-contact laser Doppler vibrometer (LDV) which can measure point-wise out-of-plane displacement responses is employed simultaneously to compare measurement results obtained from the proposed FBG displacement senor. A traditional FBG strain sensor and a strain gauge are also used to verify the performance of the FBG filter before performing the experiments for the FBG displacement sensor. The sensing locations of the FBG displacement sensor, LDV, FBG strain sensor, and strain gauge are shown in Figure 6.
Experimental Results and Discussion
The performance of the FBG filter is first demonstrated by the strain sensors (i.e., FBG strain sensor and strain gauge). A steel ball with 3.8 mm in diameter is dropped on the centerline and 10 mm away from the free end to induce transient wave propagation in the cantilever beam. Figure 7 demonstrates the measured results of the transient strain responses obtained by the traditional FBG strain sensor and strain gauge simultaneously. To observe the strain responses in detail, Figure 7 is replotted within 0.03 s and the results are shown in Figure 8. Since the transient strain responses obtained by the two strain sensors agree well with each other in Figure 8, the FBG is proved to be capable of measuring dynamic strain with high resolution both in time and space. From Figures 7 and 8, we can see that the proposed FBG filter offers excellent dynamic demodulation performance. The frequency spectrum of the smart cantilever beam can be constructed from taking the fast Fourier transform of the transient strain responses obtained by the FBG strain sensor and the results are shown in Figure 9. Next, we excite the piezoelectric actuator with random inputs to perform system identification by the stochastic spectral estimation. The random signal input is generated by dSPACE DS1104 system with the sampling frequency set to 50 kHz. From the input and output data recorded by dSPACE DS1104 system, the frequency response function (FRF) of the system is determined from the relation [21]: Figure 11. Frequency response of the smart cantilever beam obtained by the strain gauge.
From both Figures 10 and 11, we can see that a lateral mode (i.e., 602 Hz in Figure 10 and 601.4 Hz in Figure 11) is measured by the two sensors. However, frequency response obtained by the conventional strain gauge (i.e., Figure 11) contains more measurement noises especially between the first and the second bending modes of the smart cantilever beam. The identified first three bending modes of the smart cantilever beam without the presence of the steel ball are 76 Hz, 400 Hz, and 975 Hz, respectively. Since the error between the first resonant frequency is quite small, it is acceptable to say that the resonant frequencies can be obtained from the impact-induced transient responses. However, for the purpose of active vibration control, the first bending mode of the smart cantilever beam is considered to be 76 Hz in this work.
We further investigate the dynamic demodulation ability of the FBG filter and the response time for the piezoceramic actuator to achieve the steady state vibration condition. The smart cantilever beam is excited by the piezoelectric actuator at the first resonant frequency (76 Hz) and the second resonant frequency (400 Hz). Figure 12 demonstrates the dynamic strain from transient response to steady state of the smart cantilever beam excited at 76 Hz. From Figure 12, we can see that there is an overshoot in the transient strain responses before reaching the steady state. The response time from the beginning of excitation to the steady state is about 1.26 s. The transient strain responses of the smart cantilever beam excited at the second resonant frequency (i.e., 400 Hz) is demonstrated in Figure 13. The response time in this case is about 0.33 s. The same phenomenon of overshoot, although is smaller, can still be observed in the transient strain responses excited at 400 Hz. Figure 13. Transient strain responses of the smart cantilever beam excited by the piezoelectric actuator at 400 Hz. Figure 14 shows the sensing results obtained from two strain sensors simultaneously when the smart cantilever beam vibrates freely after excitation at 76 Hz is removed. Figure 14. Free vibration of the smart cantilever beam excited at first resonant frequency (76 Hz) obtained from FBG strain sensor and strain gauge. Figure 15 overlaps the two results and focuses on the responses with a span of 0.2 s. Since the responses measured from these two strain sensors correspond excellently with each other, we can again see that the FBG filter is capable of demodulating responses obtained from the FBG sensor dynamically. Finally, to see the demodulation effect of the FBG filter in an active vibration control system, the negative velocity feedback controller is utilized. As shown in Figure 16, with response before control and the control signal as comparisons, the control result obtained from the FBG strain sensor approaches to zero at about 0.3 s. Since the control effect shown in Figure 16 agrees well with the prediction that the velocity feedback can add damping of the cantilever beam, the dynamic demodulation performance of the FBG filter in an active vibration control system is demonstrated. Now we can focus on the dynamic sensing performance of the proposed out-of-plane FBG displacement sensor. First, we perform the system identification of the smart cantilever beam again by the stochastic spectral estimation with the FBG displacement sensor and a non-contact laser Doppler vibrometer (LDV). Since the LDV has high sensitivity and resolution, its measurement results are used for comparisons with the results obtained by the FBG displacement sensor. The obtained frequency responses obtained by the two sensors are shown in Figure 17. A fourteenth model obtained from the concept of constructing the Bode plot is obtained and represented as dashed lines, as shown in Figure 17.
Thus, we can see that the proposed FBG displacement sensing system can be utilized to perform system identification for the smart cantilever beam. In Figure 17, the frequency responses obtained by the FBG and LDV agree well with each other except for the low frequencies below the first bending mode. In fact, the discrepancies for low frequencies are due to the length of the random inputs recorded in the computer. To see this effect, we excite the piezoelectric actuator again with shorter random inputs (only last for 1 s), the results obtained by the two displacement sensors are represented in Figure 18 as dashed lines. Compared with the identified frequency responses (i.e., solid lines shown in Figure 18) with longer random inputs (last for 10 s), we can see that identified modes are almost the same except for the low frequencies. Repeatability of our experimental setup and the proposed FBG displacement sensor can be demonstrated in Figure 18 for frequency contents at high frequencies. Note that the lateral mode measured in strain frequency responses shown in Figures 10 and 11 is not obtained from the displacement frequency responses (i.e., Figures 17 and 18) due to the fact that the proposed FBG displacement sensor as well as the LDV is only sensitive to out-of-plane motions.
In this paper, only the first two bending modes of the cantilever beam are considered and controlled by velocity feedback control. Similar to strain experiments, we also investigate dynamic sensing performances of the proposed FBG displacement sensor by measuring transient responses of the smart cantilever beam after being excited at the first two bending modes. Figures 19 and 20 are the measurement results when the smart cantilever beam is excited at 76 Hz and 400 Hz, respectively. Figure 21 demonstrates the experimental results obtained by the two displacement sensors and simulation result obtained from the identified model of the smart cantilever beam. Compared with LDV in Figures 19 to 21, we can see that the proposed FBG displacement sensor has excellent dynamic sensing performance. From Figure 21 we can also see that the proposed FBG displacement sensor is capable of being employed to perform system identification for the flexible structures. Figure 22 and focuses on the responses within a span of 0.2 s, showing excellent agreement between two displacement sensors. The maximum peak to peak value of LDV is 2,506.44 nm. Thus, the sensitivity of the proposed FBG displacement sensor is 0.321 mV/nm by calibration with LDV. To see the feasibility of employing the proposed FBG displacement sensor in smart structures, the velocity feedback controller is used and the control results are shown in Figure 24. It is seen that the vibrations are damped out very quickly. The settling time for the free vibration to reduce to 5% of the disturbance level is 0.15 s. The agreement of the responses shown in Figure 24 between two sensors after the controller is applied demonstrates that the proposed out-of-plane point-wise FBG displacement is capable of being integrated into a smart structure to suppress the vibrations. According to vibration theory, damping ratio of the structure can be estimated from the log decrement: FBG displacement sensor before and after the velocity feedback controller is applied to the smart cantilever beam, with the control signal as a comparison. Figure 25. Active vibration control of the smart cantilever beam excited at first resonant frequency (76 Hz) obtained from FBG displacement sensor.
Next, the delay controller is applied to the system as a comparison to the velocity feedback controller. Figure 26 shows the measurement result under 3/4 phase delay control. We can see that the results obtained from differential control (i.e., velocity feedback) and delay control are almost the same. The resonant frequency of the second bending mode of the smart cantilever beam is 400 Hz. Free vibrations at 400 Hz obtained from FBG displacement sensor and LDV are shown in Figure 27 after 0 t = . The control effect of the velocity feedback controller obtained from the FBG displacement sensor is shown in Figure 28. Damping ratios of the cantilever beam obtained from FBG displacement sensor before and after velocity feedback control is applied are 0.0134 and 0.0018, respectively. Similarly, the 3/4 phase delay controller is again applied to the smart cantilever beam and the comparison of two controllers (i.e., differential controller and delay controller) is shown in Figure 29.
Conclusions
In this study, we investigate the feasibility of utilizing a fiber Bragg grating (FBG) displacement sensor to perform active vibration suppression of a smart cantilever beam. The set-up method proposed for the FBG displacement sensor allows the FBG to detect and feed back point-wise out-of-plane displacement responses. Before performing experiments, an optical full-field measurement technique called amplitude-fluctuation electronic speckle pattern interferometry (AF-ESPI) is used to provide full-field vibration mode shapes and resonant frequencies. Furthermore, an FBG filter-based demodulation technique is adopted to obtain high SNR and dynamic sensitivity and its demodulation performance is demonstrated by a traditional FBG strain sensor and strain gauge. Then, measurement ability of the proposed FBG displacement sensor is demonstrated by excellent agreements between experimental results obtained from the FBG displacement sensor and a laser Doppler vibrometer (LDV) sensor simultaneously. Both FBG strain and displacement sensors are utilized to perform system identification of the smart cantilever beam. Finally, a simple but effective velocity feedback control algorithm is used to verify the sensing performance of the proposed FBG displacement sensing system in an active structural control system. To our knowledge, this is the first time that a point-wise FBG displacement sensor has been integrated into a smart structure for performing active vibration control. | 6,532.4 | 2011-12-13T00:00:00.000 | [
"Engineering",
"Physics"
] |
Knockdown of ubiquitin-conjugating enzyme E2T (UBE2T) suppresses lung adenocarcinoma progression via targeting fibulin-5 (FBLN5)
ABSTRACT Lung adenocarcinoma (LUAD) is the main histological type of lung cancer, which is the leading cause of cancer-related deaths. Accumulating evidence has displayed that UBE2T is related to tumor progression. However, its role in LUAD has not been fully elucidated. The expression of UBE2T was detected in LUAD tissues by qRT-PCR, western blotting, and immunohistochemistry. UBE2T shRNAs were transfected into LUAD cells to analyze the consequent alteration in function through CCK-8 assay, Edu assay, transwell assay, and TUNEL staining. The potential mechanism of UBE2T was analyzed through GEPIA and verified using ChIP, EMSA, and GST pull-down assays. Furthermore, a xenograft mouse model was used to assess UBE2T function in vivo. Results showed that UBE2T level was significantly elevated in LUAD tissues and high UBE2T expression was associated with poor overall survival and disease-free survival. Results from the loss-of-function experiments in vitro showed that UBE2T modulated LUAD cell proliferation, migration, invasion, and apoptosis. The mechanism analysis demonstrated that silence of UBE2T increased FBLN5 expression and inhibited the activation of p-ERK, p-GSK3β, and β-catenin. Moreover, following knockdown of UBE2T, the cell proliferation, migration, and invasion were decreased, and sh-FBLN5 partially reverse the decrease. In in vivo experiments, it was found that UBE2T knockdown inhibits the tumor growth in LUAD. Immunohistochemically, there was a reduction in Ki67 and an increase in FBLN5 in UBE2T shRNA-treated tumor tissues. In conclusion, UBE2T might be a potential biomarker of LUAD, and targeting the UBE2T/FBLN5 axis might be a novel treatment strategy for LUAD.
Introduction
Lung cancer is the leading cause of mortality among all malignancies worldwide [1], of which lung adenocarcinoma (LUAD) is the most common subtype, accounting for 40-50% of lung cancer [2,3]. Although achievements have been made in new therapies for LUAD such as chemotherapy, immunotherapy, and molecular targeted therapy, the prognosis of LUAD is still very poor, with a 5-year survival rate of only 15% [4,5]. Therefore, further study of the molecular mechanisms of LUAD is essential for the development of new therapies against LUAD.
Ubiquitin-proteasome pathway (UPP) plays an important role in plant growth regulation, animal reproductive development, tumorigenesis and neurological diseases. E1, E2 and E3 enzymes are involved in ubiquitination progression, of which E2 plays a very important role [6,7]. Previous studies have demonstrated that the E2 enzyme ubiquitinconjugating enzyme E2D3 (UBE2D3) is involved in the regulation of cancer radiation resistance [8,9]. Another study has shown that ubiquitinconjugating enzyme E2C (UBE2C) is highly expressed in many tumors and inhibition of UBE2C inhibits tumor progression [10]. Ubiquitinconjugating enzyme E2T (UBE2T), a member of the E2 family, is located on human chromosome 1q32.1 and has a characteristic conserved domain with a size of about 16-18 kDa [11]. According to previous reports, UBE2T was used as an important member of the Fanconi signaling pathway to participate in DNA damage repair [12]. Recent studies have discovered that UBE2T is significantly increased in hepatocellular carcinoma, gallbladder cancer, and gastric cancer, and its high expression is closely associated with the tumor size, metastasis, and poor prognosis, suggesting that UBE2T may have the potential to promote the proliferation, invasion, and metastasis of malignant tumors [13][14][15]. For instance, Ueki et al. displayed that UBE2T promoted the progression of breast cancer by degrading BRCA1 [16]. Moreover, Wang et al. clarified that UBE2T down-regulation suppressed osteosarcoma cell proliferation and metastasis via inhibiting the PI3K/Akt signaling pathway [17]. More importantly, increasing evidence has shown that UBE2T is closely related to non-small cell lung cancer progression [18,19]. However, the correlation between UBE2T and LUAD, especially with regard to proliferation, invasion, migration, and apoptosis, has not yet been defined.
Fibulin-5 (FBLN5), a 66-kDa secreted glycoprotein, is identified by two independent groups in 1999 [20]. It has been reported that FBLN5 plays an important role in cell adhesion and motility, cell growth, cell metastasis, and tumorigenesis [21][22][23]. There is increasing evidence that FBLN5 has prognostic potential as a tumor suppressor in a variety of cancers, such as ovarian cancer, breast cancer, and hepatocellular carcinoma [24][25][26]. In lung cancer, overexpression of FBLN5 suppressed cell invasion and metastasis through the ERK pathway [27]. Another report demonstrated that FBLN5 impedes Wnt/ β-catenin signaling by inhibiting ERK activation of GSK3β in lung cancer [28]. Furthermore, a previous report showed that ERK/GSK3β pathway regulates cell proliferation and metastasis and is frequently activated in tumor tissues including LUAD [29]. Besides, UBE2T was reported to promote the activation of GSK3β pathway in nasopharyngeal carcinoma [30]. However, whether FBLN5/ERK/GSK3β pathway was affected by UBE2T in LUAD is not clarified.
In this study, we hypothesized that UBE2T knockdown exerted an inhibitory effect on the progression of LUAD through regulating FBLN5 expression and ERK/GSK3β signaling pathway. The purpose of this study was to elucidate the functional role of UBE2T in the proliferation and metastasis of LUAD in vitro and in vivo, and explore the molecular mechanism underlying its role.
Clinical samples
Sixty-five pairs of primary LUAD tissues and adjacent normal tissues were collected from patients admitted to the People's Hospital of Shanxi Province, after receiving written informed consent. This study was approved by the Ethics Committee of People's Hospital of Shanxi Province (approval no. SXSRM2018061325YY). This study was conducted following the ethical standards of our hospital and the Helsinki Declaration. The main clinical characteristics of the patients are summarized in Table 1.
Online database analysis
The analysis of UBE2T and FBLN5 levels in LUAD tissues from TCGA dataset was performed by using Gene Expression Profiling Interactive Analysis (GEPIA, http://gepia.can cer-pku.cn) [31]. The GEPIA and PrognoScan database (http://www.abren.net/PrognoScan) were applied to investigate the relationship between UBE2T and the prognosis of LUAD patients.
Cell lines and cell culture
The lung adenocarcinoma cell lines (H1975 and H1650) and normal human bronchial epithelial cell line (HBE) used in this study were purchased from the Cell Bank of the Chinese Academy of Sciences (Shanghai, China). All the cells were cultured in DMEM with 10% FBS and maintained in an incubator at 37°C under 5% CO 2 conditions [32].
RNA extraction and quantitative real-time PCR (qRT-PCR)
Total RNA was extracted using RNAiso reagent (TaKaRa, Dalian, China) and reverse-transcribed to cDNA using PrimeScript RT Master Mix Kit (Takara) following the manufacturer's instructions [34]. The reaction system was configured following the protocol of the SYBR Premix Ex TaqTM II kit (Thermo Fisher Scientific, Inc.), and the RNA transcript levels were performed using the Bio-rad CFX96 real-time PCR system (Biorad, USA). β-Actin was used as the internal control, and relative expression levels were calculated by the 2− ΔΔCt method. Primer sequence was shown as following: [32].
Edu immunoflurescence assay
For this, 4 × 10 3 cells/well LUAD cells were plated in a 96-well plate. After 24 h, the medium containing 50 mM Edu (100 mL) was added and incubated for 2 h at 37°C. Subsequently, the cells were fixed with 4% paraformaldehyde for 30 min and counterstained with Hoechst 33,342 for 10 min to stain the nucleus. The Edu-positive cells were counted under a fluorescence microscope (AF6000, Leica, Wetzlar, Germany) [35].
Transwell assay
Stable expression LUAD cells (1 × 10 5 cells/well) in serum-free media were placed into the upper chamber of an insert for migration assays (8-μm pore size, Corning, NY, USA) and invasion assays with Matrigel (Sigma-Aldrich, USA). The lower chambers were filled with complete medium supplemented with 20% FBS. After incubation for 48 h, the migrated or invaded cells were fixed with methanol for 10 min and stained with 0.1% crystal violet [32]. The cells were counted under a BX53 microscope in five randomly selected fields (Olympus, Japan) (magnification; 200×).
TUNEL staining
One Step TUNEL Apoptosis Assay Kit (Beyotime) was applied to measure cell apoptosis according to the manufactures protocol [36]. Briefly, LUAD cells were incubated with PBS containing 0.3% Triton X-100 for 10 min and then incubated with 0.3% H 2 O 2 in PBS for 20 min. Subsequently, the cells were incubated with TUNEL detection solution (50 μL) in the dark at 37°C for 1 h, followed by the streptavidin -HRP working solution for 30 min. Next, the cells were stained with Hoechst 33,342 for 10 min, and photographed by a fluorescence microscope (Fluoview FV1000, Tokyo, Japan) and counted by a Nikon ECLIPSE Ti fluorescence microscope under five random fields. Cell apoptosis (%) was calculated by the percentage of TUNEL-positive cells in the total number of cells (DAPI-positive cells).
Chromatin immunoprecipitation (ChIP) assay
SimpleChIP® Plus Enzymatic Chromatin IP Kit (Magnetic Beads, CA, USA) was used to perform ChIP assay according to the manufacturer's introduction [35]. In brief, LUAD cells were immobilized with formaldehyde, and the chromatin was fragmented by enzymatic hydrolysis and ultrasonic treatment. Subsequently, the chromatin was immunoprecipitated with specific anti-UBE2T, anti-FLBN5, and normal IgG antibodies. The enrichment of specific DNA fragments was analyzed by qRT-PCR.
Electrophoretic mobility shift assay (EMSA)
Biotin end-labeled probes were prepared by Sangon Biotechnology Co., LTD. (China). LightShift® Chemiluminescent EMSA Kit (Pierce, USA) was carried out to perform EMSA following the manufacturer's instructions [37]. DNA binding reactions were performed with or without anti-UBE2T antibody. DNA-protein complexes were separated by electrophoresis and then transferred onto a positive-charged nylon membrane (Millipore, USA), followed by UV light crosslinking. The signal was visualized with chemiluminescent substrate followed by film exposure.
Glutathione S-transferase (GST) pull-down assay
GST-UBE2T protein or GST control was transformed into E. coli BL21, and then 1 mmol/L IPTG was added to induce the protein expression [38]. Flag-FBLN5 protein was extracted from H1975 cell lysates. After GST-UBE2T and GST beads were cultured for 3 h, the eukaryotic expression protein Flag-FBLN5 was added to the mixture and the column was rotated vertically on the mixer for 3 h. Subsequently, the protein mixture was washed five times and denatured in 2× loading buffer at 95°C for 5 min. Protein bands were detected by western blotting.
Tumor growth in vivo
A total of 20 female BALB/c nude mice (6-8 weeks old) were purchased from the Animal Center of the Chinese Academy of Science (Beijing, China). All animal experiments were approved by the Animal Care and Use Committee of the above hospital (approval no. SXSRM2019010264YY). H1975 cells stably expressing sh-UBE2T1 or negative control (sh-NC) were subcutaneously injected into the right flank of mice (2 × 10 7 cells/mL, 0.2 mL, n = 5 in each group). The tumor volumes were recorded every 7 days (volume = length × width 2 )/2. Animals were euthanized via 2% pentobarbital sodium (120 mg/kg bodyweight) at 28 days, and then the tumors were peeled carefully and weighted [35].
Immunohistochemical (IHC) analysis
The slides (4 μm-thick) of paraffin-embedded xenograft tissues were placed in xylene and ethanol for hydration treatment. After washing 3 times with PBS, the slides were completely immersed in 95°C antigen retrieval solution for 10 min. Then, 3% H 2 O 2 were added and incubated for 10 min. After blocking with 5% serum for 30 min at 37°C, the slides were probed with specific rabbit anti-UBE2T antibody, mouse anti-FBLN5 antibody and rabbit anti-Ki67 at 4°C overnight, followed by the HRP goat anti-rabbit/mouse IgG for 1 h. Thereafter, DAB dropwise to the slides, and the color reaction was terminated with tap water. Finally, the slides were photographed by Olympus BX40 microscope (Tokyo, Japan) and quantified with Image ProPlus (IPP) software (Media Cybernetics, Rockville, MD, USA) [32].
Statistical analysis
Statistical analysis was performed by using GraphPad Prism 8.0.2. Two-tailed unpaired Student's t test and Turkey's post hoc tests in oneway ANOVA were applied to analyze the differences between two groups or among the multiple groups. The χ 2 test was used to evaluate the association between the expression of UBE2T1 and clinicopathological characteristics of LUAD patients. All data are shown as the mean ± SD. P < 0.05 was considered as a statistically significant.
Results
This study was aimed to elucidate the functional role of UBE2T in the proliferation and metastasis of LUAD in vitro and in vivo and to explore the molecular mechanism underlying its role. Results showed that UBE2T was significantly elevated in LUAD tissues and high UBE2T expression was associated with poor overall survival. Results from the loss-of-function experiments in vitro showed that UBE2T modulated LUAD cell proliferation, migration, invasion, and apoptosis. The mechanism analysis demonstrated that silence of UBE2T increased FBLN5 expression and inhibited the activation of ERK/GSK3β pathway. In in vivo experiments, it was found that UBE2T knockdown inhibits the tumor growth in accumulating evidence has shown that that FBLN5, as a tumor suppressor, has LUAD. These results suggested that UBE2T might be a potential biomarker of LUAD, and targeting the UBE2T/FBLN5 axis might be a novel treatment strategy for LUAD.
High expression of UBE2T was closely related to the poor prognosis of LUAD patients
To determine the expression status of UBE2T in LUAD, the online tool GEPIA was applied to analyze TCGA LUAD dataset. Results revealed that the UBE2T mRNA expression was elevated markedly in LUAD tissues versus to that in normal tissues (Figure 1(a)). In addition, UBE2T expression was evaluated in 65 pairs of LUAD tissues. In line with TCGA dataset results, UBE2T was significantly upregulated in LUAD tissues in comparison with normal tissues (Figure 1(b)). Similarly, the protein level of UBE2T was increased in LUAD tissues (Figure 1(c)). Next, the relationship between UBE2T and the survival rates of LUAD patients was evaluated. According to the median level of UBE2T, patients were divided into lowexpression group and high-expression group. Results displayed that high UBE2T expression was closely associated with the poor overall survival and poor disease-free survival in LUAD patients (Figure 1(d,e)). Another online database, PrognoScan database, was then used to examine the prognostic potential of UBE2T in LUAD. As shown in Figure 1(f,g), there was a poor prognosis in LUAD patients with high UBE2T expression than in those with low UBE2T expression. Therefore, these findings suggested that UBE2T might be a potential prognostic biomarker for LUAD patients.
UBE2T knockdown suppressed LUAD progression
To investigate the functional roles of UBE2T in LUAD, UBE2T expression was first measured in LUAD cells. Results from qRT-PCR showed that UBE2T was higher in LUAD cells than that in normal human HBE cells (Figure 2(a)). H1975 and H1650 cells were infected with sh-UBE2T and sh-NC. qRT-PCR results showed that UBE2T was down-regulated in sh-UBE2T LUAD cells versus to that in control group (Figure 2(b)), sh-UBE2T-1 was selected for further experiments for its stronger inhibitory effects. CCK-8 assays displayed that the proliferation rate of H1975 and H1650 cells treated with sh-UBE2T-1 was slower than that with sh-NC (Figure 2(c)). Edu assays showed that UBE2T knockdown decreased the percentage of Edu-positive cells in comparison with control group (Figure 2(d)). Transwell migration assay discovered that H1975 and H1650 cells with sh-UBE2T-1 had less migrated cells than that of cells with sh-NC (Figure 2(e)). Moreover, results from transwell invasion assay revealed that H1975 and H1650 cells with sh-UBE2T-1 had less invasive cells than that of cells with sh-NC (Figure 2(f)). Besides which, western blotting assay discovered that Bcl-2 level was downregulated, while Bax and cleaved-caspase-3 levels were upregulated in sh-UBE2T-1 LUAD cells, relative to sh-NC group (Figure 2(g)).
Meanwhile, TUNEL staining revealed that the apoptosis cells was obviously increased in sh-UBE2T-1 group (Figure 2(h)). These data indicated that UBE2T knockdown inhibited LUAD cell proliferation, migration, invasion, and promoted cell apoptosis.
UBE2T knockdown contributes to the increase of FBLN5 and inactivation of ERK/GSK3β pathway
To investigate the underlying mechanism of UBE2T in LUAD, cBioPortal database (https:// www.cbioportal.org) was applied to search for UBE2T-related genes. Among those genes, FBLN5 was a negatively correlated gene of UBE2T in LUAD r = −0.54, p = 0.0124, Figure 3(a). Besides, GEPIA with the Spearman correlation test showed a negative correlation of UBE2T with FBLN5 mRNA expression r = −0.45, p < 0.001, Figure 3(b). In addition, the relationship between UBE2T and FBLN5 expression was evaluated in LUAD tissues as well. In line with the online database results, UBE2T was negatively correlated with FBLN5 in LUAD tissues (Figure 3(c)). Furthermore, the effect of UBE2T on FBLN5 activity was investigated in LUAD. qRT-PCR results displayed that UBE2T knockdown increased FBLN5 expression in LUAD cells (Figure 3(d)). Similarly, results from western blotting revealed that UBE2T down-regulation increased the level of FBLN5 (Figure 3(e)). Importantly, ChIP-qPCR assays displayed that UBE2T enhanced the enrichment of FBLN5 in LUAD cells (Figure 3(f)). A subsequent EMSA discovered that when the native probe was incubated with purified UBE2T protein, a DNA-protein complex was formed, and the mutated probe reduced the binding capacity. An antibody supershift assay further verified that UBE2T directly bound to native probe (Figure 3(g)). In addition, results from GST pull-down assay demonstrated that there was a direct interaction between GST-UBE2T and Flag-FBLN5 (Figure 3(h)). These results indicated that FBLN5 might be a potential target gene of UBE2T in LUAD. Next, the expression status of FBLN5 was determined in LUAD, TCGA database results displayed that FBLN5 was significantly decreased in LUAD tissues, relative to that in normal tissues (Figure 3(i)). Furthermore, in our own cohort, FBLN5 was obviously reduced in LUAD tissues, when compared to normal tissues (Figure 3(j)). Western blotting results discovered that compared with sh-NC LUAD cells, the levels of p-ERK, p-GSK3β, and βcatenin were down-regulated in LUAD cells with sh-UBE2T-1 (Figure 3(k)). These findings suggested that UBE2T knockdown suppressed the activation of ERK/GSK3β pathway in LUAD.
FBLN5 knockdown abrogated the inhibitory effect of sh-UBE2T on LUAD progression
Given that UBE2T modulated LUAD cell proliferation and metastasis and regulates FBLN5 expression in LUAD cells, the effect of FBLN5 on UBE2Tregulated cell proliferation and metastasis was further explored. H1975 and H1650 cells were transiently transfected with sh-NC, sh-UBE2T-1, or sh-UBE2T-1+ sh-FBLN5. Results from CCK-8 showed that H1975 and H1650 cells with sh-UBE2T-1 exhibited a decrease in the cell viability and FBLN5 depletion partially reversed sh-UBE2T-inhibited cell viability (Figure 4(a)). Meanwhile, Edu assay results discovered that the positive cells generated in sh-UBE2T-1+ sh-FBLN5 groups were notably enhanced versus to that in sh-UBE2T-1 group (Figure 4(b)). Moreover, results of Figure 4(c) show that the decrease in the migration resulting from UBE2T knockdown was partially rescued by silence of FBLN5. Similarly, as shown in Figure 4(d), there was a significant decrease in invasion abilities in UBE2T knockdown cells, and the decrease was restored by combining with FBLN5 inhibition (Figure 4(d)). In addition, western blotting revealed that sh-UBE2T-1 remarkably reduced Bcl-2 level and accelerated Bax and cleaved-caspase-3 levels in LUAD cells, which were partially overturned by combining with FBLN5 depletion (Figure 4(e)). TUNEL assay manifested that the promotion effect of sh-UBE2T-1 on LUAD cell apoptosis was also partly attenuated by the combined with sh-FBLN5 (Figure 4(f)). Collectively, these results confirm that knockdown of UBE2T suppressed the proliferation and metastasis potential possibly by activating FBLN5 expression in LUAD cells.
UBE2T knockdown inhibited tumor growth in vivo
To determine the effect of UBE2T on LUAD in vivo, H1975 cells with sh-UBE2T-1 were injected into nude mice to establish xenograft tumor models. Twenty-eight days after inoculation, the mice were euthanized and dissected ( Figure 5(a)). The tumor volume and tumor weight were remarkably downregulated in the sh-UBE2T-1 group compared to sh-NC group (Figure 5(b,c)). IHC results showed that the tumor cell proliferation marker Ki67 was decreased in the sh-UBE2T-1 group, relative to the control group ( Figure 5(d)). Moreover, UBE2T was obviously decreased, while FBLN5 was increased in the tumor tissues of sh-UBE2T-1 group compared with sh-NC group ( Figure 5(e,f)). These results indicated that UBE2T knockdown inhibited LUAD tumor growth in vivo.
Discussion
The dysregulation of UBE2T is closely related to the occurrence and development of various tumors [39,40]. Although a previous study has shown that UBE2T was elevated markedly in non-small cell lung cancer tissues and ranked first according to the hazard ratio in the survival analysis [18], the functional role and molecular mechanism of UBE2T in LUAD proliferation and metastasis remain unknown. Herein, we discovered that UBE2T was upregulated in LUAD tissues in TCGA dataset, as well as in LUAD tissues and cells. The results were similar to the previous study that UBE2T was overexpressed in lung cancer, which was confirmed by western blotting, qRT-PCR, and immunohistochemistry [41]. Notably, analyzing TCGA survival data and PrognoScan database demonstrated that high UBE2T expression strongly indicated poor prognosis of LUAD patients. Collectively, these findings indicated that UBE2T might be a promising biomarker for the prognosis and diagnosis of LUAD patients. Moreover, the loss-of-function assay results revealed that UBE2T knockdown suppressed LUAD cell proliferation, migration, and invasion in vitro and tumor growth in vivo. Taken together, these findings revealed that UBE2T exhibited the critical oncogenic roles in LUAD and UBE2T was indicated as a potential therapeutic target for LUAD. ChIP-qPCR assay, EMSA assay, GST pull-down assay, and TCGA data analysis were applied to explore the mechanisms of UBE2T in LUAD. Through cBioPortal database and TCGA data analysis, the potential targets of UBE2T were screened. Among those genes, FBLN5 may be a possible factor involved in the development of LUAD. FBLN5, an extracellular matrix protein, takes part in regulating the proliferation, invasion and angiogenesis of malignant tumor cells [2][3]. It has been reported that FBLN5 was significantly down-regulated in ovarian cancer [24], breast cancer [42], and lung cancer [28], and FBLN5 overexpression inhibited cells proliferation and metastasis. In line with previous results, the analysis of online databases and our findings discovered that FBLN5 was down-regulated in LUAD tissues, and statistical analysis showed a negative relationship between UBE2T and FBLN5 expression. Our findings also displayed that UBE2T knockdown activated FBLN5 expression by qRT-PCR, western blotting. Moreover, UBE2T bound to FBLN5, which was confirmed by ChIP-qPCR, EMSA and GST pull-down assays. The FBLN5 up-regulation might be crucial for UBE2T-mediated cell proliferation, migration and invasion. Thus, the effects of FBLN5 on UBE2T-regulated cell proliferation and metastasis was investigated. The rescue experiments revealed that the inhibitory effect of sh-UBE2T on the proliferation, migration, and invasion were reversed by sh-FBLN5.
Given previous studies shown that FBLN5 was involved in lung cancer development via the ERK pathway. A recent study showed that p-ERK promotes β-catenin activity by suppressing its regulatory molecule GSK3β [43]. Our study demonstrated that silence of UBE2T decreased the levels of p-ERK, p-GSK3β, and βcatenin, indicating that UBE2T knockdown impeded the activation of ERK/GSK3β signaling pathway.
There are limitations in this present study. First, only H1975 cells used in in vivo experiments, due to the limitation of time and technology. Second, the mechanisms of LUAD are complex and the targets of UBE2T are diverse; several crosstalk signaling pathways participate in the network regulated by UBE2T in LUAD, it is necessary to further study how UBE2T regulates FBLN5 and ERK/GSK3β signaling pathway. Third, there are few clinical research in this paper, which should be performed in further studies.
Conclusions
This present study showed that UBE2T deficiency suppressed LUAD progression through increasing FBLN5, and suppressing the activation of ERK/ GSK3β pathway. These results elucidate the underlying mechanism by which UBE2T regulates LUAD progression and provides a new direction for the development of effective treatment strategies for LUAD.
Disclosure statement
No potential conflict of interest was reported by the author(s).
Funding
The author(s) reported there is no funding associated with the work featured in this article.
Availability of data and materials
The datasets used and/or analyzed during the present study are available from the corresponding author on reasonable request.
Authors' contributions
Yi Li designed, supervised the research, and participated in writing the manuscript. Xiaojuan Yang performed the literature researches and the experimental studies. Dan Lu performed data analysis and statistical analysis. All authors read and approved the final manuscript.
Ethics approval
The study was approved by the Ethics Committee of People's Hospital of Shanxi Province (approval no. SXSRM2018061325YY).
Consent for participate
Signed written informed consents were obtained from the patients and/or guardians. | 5,525.8 | 2022-05-01T00:00:00.000 | [
"Medicine",
"Biology"
] |
Nadir ozone profile retrieval from SCIAMACHY and its application to the Antarctic ozone hole in the period 2003-2011
The depletion of the Antarctic ozone layer and its changing vertical distribution has been monitored closely by satellites in the past decades ever since the Antarctic "ozone hole" was discovered in the 1980’s. Ozone profile retrieval from nadir-viewing satellites operating in the ultraviolet-visible range requires accurate calibration of level-1 (L1) radiance data. Here we study the effects of calibration on the derived level-2 (L2) ozone profiles and apply the retrieval to the Antarctic ozone hole region. 5 We retrieve nadir ozone profiles from the SCIAMACHY instrument that flew on-board Envisat using the Ozone ProfilE Retrieval Algorithm) (OPERA) developed at KNMI with a focus on the stratospheric ozone. We study and assess the quality of these profiles and compare retrieved (L2) products from L1 SCIAMACHY versions 7 and 8 indicated as respectively (v7, v8) data from the years 2003-2011 without further radiometric correction. From validation of the profiles against ozone sonde measurements, we find that the v8 performs better due to correction for the scan-angle dependency of the instrument’s optical 10 degradation. The instrument spectral response function can still be improved for the L1 v8 data with a shift and squeeze. We find that the contribution from this improvement is a few percent residue reduction compared to reference solar irradiance spectra. Validation for the years 2003 and 2009 with ozone sondes shows deviations of SCIAMACHY ozone profiles of 0.8%−15% in the stratosphere and 2.5%− 100% in the troposphere, depending on the latitude and the L1 version used. Using L1 v8 for the 15 years 2003-2011 leads to deviations of∼ 1%−11% in stratospheric ozone and∼ 1%−45% in tropospheric ozone. Application of SCIAMACHY v8 data to the Antarctic ozone hole shows that most ozone is depleted in the latitude range from 70◦S to 90◦S. The minimum integrated ozone column consistently occurs around 15 September for the years 2003-2011. Furthermore from the ozone profiles for all these years we observe that the value of the ozone column per layer reduces to almost zero at a pressure of 100 hPa in the latitude range of 70◦S to 90◦S, as was found from other observations. 20
Introduction
Ozone (O 3 ) is one of the most important trace gases in our atmosphere.Stratospheric O 3 absorbs the dangerous solar ultraviolet (UV) radiation making it an important protector of life.A small amount of O 3 is found in the troposphere originating from air pollution and photochemistry -this ozone is considered a health-risk.Daily ozone monitoring using satellites dates back to the late 1970s with instruments like the Total Ozone Monitoring Spectrometer (TOMS, 1979) and Solar Backscatter Ultra Violet (SBUV) instruments, and since the mid-1990s also by the full UV/VIS spectrum covering satellite instruments like the Global Ozone Monitoring Experiment (GOME, GOME-2) (e.g.Burrows et al., 1999;Munro et al., 2016), Scanning Imaging Absorption spectroMeter for Atmospheric ChartograpHY (SCIAMACHY) (Bovensmann et al., 1999), and Ozone Monitoring Instrument (OMI) (Levelt et al., 2006), to name a few.These successions of instruments allow us to compare long term global ozone layer behaviour and cross-check the quality of the measured data.
A long-term monitoring of the ozone trend layer is primarily driven by global measurements of total ozone time series of such satellite data.
The vertical profile of ozone has traditionally been measured by in-situ electrochemical instruments attached to balloons (ozone sondes).Although ozone sondes provide the most accurate method for ozone measurement, they are limited in the heights they can reach (< 35 km), and their geographical coverage is limited to approximately 300 stations worldwide that provide weekly ozone sonde profiles, and very few stations with a higher than weekly measurement frequency.Satellite measurements provide an alternative means for obtaining globally vertical ozone profiles.In general limb and occultation mode satellite instruments can well resolve the vertical distribution in stratospheric ozone.However, they are limited in their horizontal resolution, and they have no sensitivity to ozone in the middle and lower troposphere.An alternative approach is to use satellite measurements in nadir mode by high-resolution spectrometers in the thermal IR, like IASI and TES, and in the UV/VIS, like GOME, GOME-2 (e.g.Cai et al., 2012;van Peet et al., 2014;Keppens et al., 2015;Miles et al., 2015).OMI (e.g.Liu et al., 2010;Kroon et al., 2011), and SCIAMACHY.
The observation principle of nadir ozone profile retrieval in the UV/VIS is based on the strong spectral variation of the ozone absorption cross-section in the UV-visible wavelength range, combined with Rayleigh scattering.The key here is that the short UV wavelengths (265-300 nm) are back-scattered from the upper part of the atmosphere whereas the longer UV wavelengths (300-330 nm) are mostly back-scattered from the lower part of the atmosphere.This transition in the ozone crosssection between 265-330 nm is useful in retrieving its vertical profile.Nadir UV and visible spectra provide better horizontal resolution in ozone although their observations can only be carried out in daytime.In the thermal infra-red measurements can be done during both night and day.SCIAMACHY had both limb and nadir mode capability.There have been several studies of ozone profiles using SCIAMACHY limb data (e.g.Brinksma et al., 2006;Mieruch et al., 2012;Hubert et al., 2017).Brinksma et al. (2006) found biases in ozone profile of stratosphere < 10%.Also in their analysis of limb scatter ozone profiles from 2002-2008, Mieruch et al. (2012) ) found the stratospheric ozone have a bias of ∼ 10% against correlative data sets and this bias increased up to 100% in the troposphere for the tropics.Similarly Hubert et al. (2017) in their more recent study of limb profiles find that the SCIAMACHY ozone biases are about ∼ 10% or more in the stratosphere with short-term variabilities of ∼ 10%.There has been very little published work on ozone profile retrieval from SCIAMACHY nadir mode, probably due to calibration issues.
In this study, we focus on nadir ozone profile retrieval from SCIAMACHY and the impact of L1 calibration improvements.
We evaluate the three most recent versions of the SCIAMACHY L1 product dataset (described in Sect.2) on the basis of retrieved nadir ozone profiles from the SCIAMACHY UV reflectance spectra.The result of this paper shows the improved quality of the latest L1 dataset version.
The results presented here highlight the need for further corrections of the L1 data.However, a detailed study of radiometric bias corrections in L1 data is beyond the scope of this paper.The focus of this study is to analyse stratospheric ozone and this study is done for almost the entire mission length of 2003-2011 where validation is done globally for the latest L1 version available (v8).We will focus exclusively on the study of ozone in the stratospheric region (100-10 hPa), but we will briefly comment on the accuracies we get for tropospheric region (1000 -100 hPa).The OPERA retrieval algorithm is briefly reviewed in Sect.2.2.In Sect.3, the analysis of the slit function of the instrument is presented.Results on the ozone profiles, and the comparison between the level-1 datasets are shown in Sect. 4. This is followed by comparison to sondes in Sect. 5 for the most recent dataset from years 2003-2011 spanning almost the entire mission.We apply the SCIAMACHY dataset in analysing the Antarctic ozone hole in Sect.6.We discuss the possible effects of L1 radiometric bias corrections and applying slit function corrections in Sect.7, and finally conclude in Sect.8.
Instrument, data and methods
SCanning Imaging Absorption spectroMeter for Atmospheric ChartograpHY (SCIAMACHY) is space-borne spectrometer on board ESA's Environmental Satellite Envisat (Burrows et al., 1995;Bovensmann et al., 1999) with both horizontal (limb) and vertical (nadir) mode viewing design covering the wavelength range from 212 nm (UV) to 2386 nm (infrared, IR) spread over 8 channels.Launched in March 2002, its mission lifetime spanned until 2012 and we have level-1 data (Lichtenberg et al., 2006) from the spectrometer from August 2002 until April 2012.In this paper we concern ourselves with the nadir data and in retrieving the vertical distribution of ozone using 265-330 nm UV-VIS continuous spectral data.Each nadir state is an area on the Earth's surface defined by the scan speed of the nadir mirror across track direction and the spacecraft speed in the along track direction, the field of view (FoV) and the operation of the instrument.This gives typical ground pixel sizes (or equivalently nadir states) of 240 km × 30 km for an integration time (IT) of 1.0 s and 60 km × 30 km for an IT of 0.25 s (Gottwald and Bovensmann, 2011).Alternatively the nadir viewing corresponds to an instantaneous FoV (IFoV) of 0.045 separates the clusters as listed in Table 1.
(across track) × 1.8 • (along track).The IT also varies for clusters, the measurements should be combined before it is fed to the optimal estimation.The variation in IT between clusters can give rise to a spectral jump (discontinuity) where the last value of the spectral value in the preceeding cluster does not match the first value of the same in the following cluster.The wavelengths at which this occurs for Channels 1 and 2 were identified and blocked from our analysis (see Fig. 11 in Sect.7).
The above specified wavelength range straddles over Channels 1 and 2 of SCIAMACHY with an overlap between the two channels.We use wavelengths 265-314 nm from Channel 1 and 314-330 nm from Channel 2. The extracted data from L1 are broken into states, which are groups of ground pixels.The spectrum of each ground pixel is divided into spectral clusters which are groups of wavelengths having their own integration time.These are then organized into clusters of data instead of channels.
The mapping of clusters to wavelengths, the resolution, and IT are listed in Table 1.
The most important input for retrieving ozone profile are the Earth's reflectance spectra (from the L1 product).An example of these spectra are shown in Fig. 1.The observed reflectance spectrum is defined as: where I, µ 0 , and E are the radiance scattered and reflected by the Earth's atmosphere, the cosine of the solar zenith angle (θ 0 ), and the incident solar irradiance at the top of the atmosphere perpendicular to the solar beam, respectively.
Versions of level-1 (L1) data
We make use of three different versions of L1 products (described below) from the nadir spectral data and present their differences using L2 products (ozone retrieval).Specifically the calibrated L1c data are reproduced using ESA's SciaL1C program of v3.2.6 1 .These versions we use in this paper are described below: 1 2. v7 mfac : This data set is identical to the one described above, except here we use the degradation corrections that were provided independently as auxiliary data files 2 .Thus, the data structure allows us to turn on and turn off to check the effects of the degradation corrections independent of other calibration corrections.Degradation correction is obtained by the so-called 'm-factors'.The m-factors are determined by the monitoring of the light path which is given by the ratio of the measured spectrum of a constant source (Sun) to that obtained for the same optical path at a given time.This gives therefore an indication of the part of the degradation of the optical path as the instrument ages.These m-factors are simple multiplication factors to the solar spectra after the absolute radiometric calibration (Gottwald and Bovensmann, 2011).
3. v8: The IPP v8.02 is the 2016 version of SCIAMACHY L1 product.The main difference between this version with the ones above is the implicit implementation of a standard degradation correction.The degradation in this version takes into account the scan angle dependence of the nadir viewing geometry of the instrument with the optical path.We use the slit-function key data provided in v8 for the instrumental slit function retrieval (see Sect. 3).Specifically the radiometric calibration uses a scan mirror model which takes into account the physical effect from the contamination layers in the mirror.The degradation using this model gives a scan angle dependence (Bramstedt, 2014).
OPERA retrieval algorithm
The Ozone Profile Retrieval Algorithm (OPERA) has been developed in KNMI (van Oss and Spurr, 2002;van der A et al., 2002).It retrieves the vertical ozone profile using nadir satellite observations of back scattered UV sunlight from the atmosphere using UV and visible wavelengths.The algorithm makes use of the laws of radiative transfer in computing the top-ofatmosphere radiances given a number of atmospheric scattering and absorption parameters.The ozone absorption cross-section decreases from 265 nm to 330 nm which allows us to retrieve the amount of ozone as a function of atmospheric height.The retrieval method is based on a forward model with a maximum posteriori approach following Rodgers (2000).This amounts to obtaining the state of the atmosphere by using the radiative transfer model and inversion technique iteratively till the model atmosphere matches the measurement.For a comprehensive algorithm overview and retrieval configuration, along with a description of the evaluation of the algorithm and the application to GOME-2 data, we refer to (Mijling et al., 2010;van Peet et al., 2014).The configuration chosen for all the retrievals are tabulated in Table 2.The retrieval grid or the vertical resolution of the nadir profile in OPERA can be chosen according to the Nyquist criterion.For SCIAMACHY data we find that setting the vertical grid to 32 layers or more gives the same value for the degrees-of-freedom (DFS).DFS (used in the Results section below) is a number related to the averaging kernels of the instrument or the sensitivity of the instrument with vertical height.
In practice, the measurement (reflectance spectrum) R meas (λ) in Eq. 1 is prepared in the beginning of the OPERA algorithm, which is then passed to the Forward model.This model contains vertical atmospheric profiles, temperature, a priori ozone profile, geolocation, cloud data and surface characteristics.The forward model is used to compute simulated radiance at the top-of-atmosphere at wavelengths determined from the measured instrumental spectral data.This is further used to generate reflectances by using convolved simulated solar irradiance spectrum.The inversion step that follows is based on the Optimal Estimation method requiring measurement, simulation and measurement uncertainties in vector/matrix forms.An inversion using derivatives of reflectances with respect to the desired parameter to be solved is carried out until convergence is reached or until the maximum number of iterations is reached.For a comprehensive description of the flow of the algorithm we refer to the OPERA manual (Tuinder et al., 2014).
Instrumental Slit Function calibration of the Solar spectral measurement
The accuracy of the retrieved geophysical product is primarily driven by the quality of the measured spectra, R meas (λ) (Eq.1), and its spectral and radiometric calibration.The most important SCIAMACHY specific calibration applied to the level-0 (raw data) is described in Slijkhuis et al. (2001).One of the spectral calibrations done often to assess the degradation of the instrument in-flight is a fit of the instrument slit function (SF).It describes the behaviour of the projection of the incoming light onto the detector pixels.These were measured for SCIAMACHY on ground before launch ( 2002) and they are provided as L1 key data for the instrument that measures the solar spectra, E(λ) (Eq.1).The slit function of the instrument has different functional forms depending on which channel they belong to.For Channels 1-2 (relevant for our ozone retrieval) they are described by a single hyperbolic function: The L1 key data provides the full width half maximum (FWHM) measured on ground which is related to the a parameter in the equation above.This parameter can be solved in terms of the FWHM by using the fact that at the central point (λ 0 ): shift, D and squeeze, S are applied to Eq. 2: followed by radiometric parameters, gain G and offset O: The unit of Shift, D is nm and of Offset, O, Gain, G and Squeeze, S are numeric factors and all four parameters depend on wavelength.The parameters FWHM and λ 0 at each wavelength taken from the v8 L1 key data were identical for all the solar spectra throughout the mission.These values are given at certain wavelengths spread throughout Channels 1 and 2 and were interpolated for the wavelengths in between.An Optimal Estimation (OE) algorithm is used to solve for the best fit parameter values using the solar measurements after the launch.The retrieved values at the solar spectrum wavelengths are then also applied to the Earth radiance spectra I(λ) (Eq. 1) interpolated at the solar spectrum wavelengths.
Thus for each solar spectrum wavelength the parameter values are used in the retrieval algorithm.These best fit values are computed using the slit function model with the above mentioned four ways of manipulation.The manipulations are that each spectral peak of the model spectrum (reference) can be transformed with: Offset, Gain, Shift, and Squeeze.Each of these spectral manipulations can be modelled as a polynomial of order n for each slit function of the instrument at the desired channels.The OE was run using this model from which the best values of the Offset, Gain, Shift, Squeeze as a function of wavelength of Channels 1 and 2 are retrieved.The best fit value is checked by evaluating the relative difference between the Solar Irradiance Measurement and Simulation which is the relative residual.The relative residuals for the best fit {O, G, D, S} are shown in Fig. 2. Each curve is a residual for one solar spectrum where the blue line is achieved by using only gain and offset in the model and the red line is achieved by including the shift and/or squeeze.
We find that in Channel 1, the wavelength range 265-308 nm is the optimal range for obtaining smallest residuals.However because the ozone profile retrieval algorithm requires data from 260 nm, we do the optimal estimation from 260-308 nm.The higher wavelength values around 314 nm in Channel 1 are known to have had calibration problems and therefore we use them only until 308nm where the slit function retrievals are still well behaved.The irradiances at 308-314 nm in Channel 1 failed to match the model spectrum.We ran the optimal estimation on all solar measurements of the SCIAMACHY mission for each day which amounted to 3463 solar spectra.We convolved all spectra with {Offset, Gain, Shift, Squeeze} of polynomial orders of = {2, 15, 2, −} and found that the cost function for the majority of cases reaches less than 1 (as expected) within 10 loops of retrieval in the Optimal Estimation routine.We find that using Shift in range 265-308 nm in Channel 1 reduces the residuals significantly looking at the left-panel in Fig. 2.
The residuals in Channel 2 are larger throughout.Dividing the relevant range of 308-330 nm into smaller ranges to get smaller residuals did not reduce the residuals any further.For the optimal wavelength divisions, we found the least residuals given by the polynomial orders of the set of {Offset, Gain, Shift, Squeeze} = {1, 4, 1, 2}.In Channel 2, shown in the right panel, we find that using Shift, Squeeze in the range 308-330 nm reduces the relative residuals significantly.The relative errors of the observation in the wavelength range which we will use for ozone retrieval, 265-330 nm, are in the order of 10 −5 .The relative residuals are in the order of a few percent.So we expect an error of a few percent to propagate into the ozone retrieval despite more accurate solar irradiances.There are anomalies in the residuals at around 279 nm, 280 nm and 285 nm.These are due to the strong MgI and MgII lines from the solar spectra and will not be used for ozone retrieval as indicated in Table 2. in red deviates from the other two visibly, which is hard to see in the reflectance spectra in the left panel.These differences are exacerbated for the year 2009 (later time of the mission) where the profile of v7 is significantly different from v7 mfac and v8 in its shape and amount of ozone.The corresponding measured spectra in the bottom-left-panel confirm these differences.
We also observe visible differences between v7 mfac and v8 indicating the intrinsic differences in the implementation of the degradation corrections between the two dataset versions.In the right panels of the figure are horizontal black-dashed lines demarcating the lower-middle stratosphere (100-10 hPa) with another line at 50 hPa.In dataset v7 there are large variations in the troposphere (1000-100 hPa) and a significant reduction of the peak of the ozone value in the stratosphere, suggesting the unreliability of this dataset for later years of the SCIAMACHY mission.The median errors and standard deviations (st. dev.) along with number of pixels and convergence statistics for the retrievals in Fig. 4 are listed in Table 3.The maximum number of iterations, n_iter (see Table 2), is set to 10 and the median values of this quantity in the table conversions.This also suggests that further corrections to the L1 data are needed (See Sect.7).
Validation: Comparison with ozone sondes
For validation, the retrieved ozone profiles are compared with balloon ozone sondes obtained from the World Ozone Ultraviolet Radiation Data Centre (WOUDC, 2011).The sonde is used if it is located within the four corners of the SCIAMACHY pixel/state on the ground and has a measurement date and time within 6 hours of the sonde.
The resulting geolocations of the selected sondes from the WOUDC data set are plotted in Fig. 5 for all years.The validation algorithm implementing the collocation criteria is very similar to the one used by van Peet et al. (2014).For the methodology we refer to that paper.The number of stations for years 2003-2011 range from 38 to 66 with number of sondes ranging from 1 to 55 per station.In comparing the retrieved profile with that of the sonde, the sonde profile is convolved with the averaging kernel (AK) from the retrieval characterizing the sensitivity of SCIAMACHY at each layer.This gives a smoothed sonde profile which is more suitable to compare with the profile retrieved from the satellite instrument and influences the results in Figs.
6-7 which are discussed below.To get an impression of the shape and effect of the averaging kernels, we refer to Appendix A. There in Figure A1 we show an example of SCIAMACHY averaging kernel shapes for a subset of the retrieved layers.In The relative percentage differences between the nadir profiles and the sondes are given in Table 4, where the top half of the table shows the results for the years 2003 and for datasets v7, v7 mfac and v8 for each zone: SH, Tr and NH.The same is listed in the bottom half for the year 2009.From the table we see that the absolute values of the deviations (see third column, st.dev.[%]) in the stratosphere are systematically smaller than in the troposphere (see fifth column).These large spreads in st.dev. of satellite retrievals from sondes are also visible in Fig. 6.Any deviation above ∼ 15% in the stratosphere and above 20% in the troposphere are shown in bold for reference as these are the required accuracy levels in the ESA CCI programme (http://www.esa-ozone-cci.org/).This is probably not due to the relatively worsening sensitivity of the nadir instrument to the tropospheric retrievals (van Peet et. al., private communication).Rather it might be due to the quality of the SCIAMACHY data.Deviations in validation in the upper troposphere and lower stratosphere have been reported due to the ozone variability in a previous study of GOME-2 nadir data (Cai et al., 2012).However, in the stratosphere (within the yellow dashed lines), the and sonde (convolved with satellite AK) normalised by the sonde profile and the dashed lines correspond to the difference with the a priori profiles used in the retrieval (see Table 2).The number of collocated pixels for all the latitude bands for each year is listed in Table 5 along with median degrees-of-freedom (DFS, see Sect.2.2), median number of iterations required to achieve convergence, the median solar azimuth angle and the median deviations in % for each zone in stratosphere and troposphere in 5 show often values higher than 20% (in boldface) for troposphere whereas these values are less than 15% for the stratosphere deeming them within specifications according to ESA CCI requirements (see Sect.
SCIAMACHY results of the Antarctic Ozone Hole
In this section we show an application of the SCIAMACHY dataset using L1 v8 to infer the ozone in the Antarctic region.
We apply the OPERA retrieval algorithm to all ground pixels south of 45 5.The median uncertainties in the retrieved columns and their st.dev.are also listed in Table 5.The variations in the minimum integrated ozone are lowest for the latitude band: [45 • S:55 • S] for all years (see top row of the figure) instead of a V-shaped dip appears in the southernmost latitude bands (see middle and bottom panels in the figure).This is as expected as the ozone depletion is stronger in the Antarctic region.Also the time of the minimum of the daily minimum integrated ozone columns occur between 15 September -15 October which is also expected from other published results.
We compare these time series with the Multi Sensor Reanalysis of Ozone, version 2 (MSR v2) (van der A et al., 2015).The MSR uses data from TOMS, SBUV, GOME, SCIAMACHY, OMI, GOME-2 for latitudes south of 30 • S resulting in a multidecadal ozone column.The reprocessed ozone columns from all satellites are assimilated and bias corrected by calibrating with included in the MSR v2 makes still use of v7 using the the total nadir ozone retrieval algorithm (TOSOMI) (Eskes et al., 2005;Valks and van Oss, 2003).The minimum ozone value available for any day from the MSR dataset is plotted for latitude band of 70 • S:90 • S in Fig. 9 and can be qualitatively compared to the bottom two panels in Fig. 8.The characteristic "V" shape in the It is useful to specifically compare the SCIAMACHY ozone total column in the stratosphere to the MSR v2 dataset for the year 2010, which was found to be an anomalous year in its behaviour of the ozone depletion (de Laat and van Weele, 2011).
The Antarctic ozone hole had 40-60% less ozone destruction compared to the average of previous years (2005)(2006)(2007)(2008)(2009).This is reflected in the right-panel of Fig. 9 where the minimum integrated ozone (in red) are above the rest of the years, showing higher levels of ozone from 1 August to 1 November with variations in the range of 250-150 DU.Comparing this with the right panel of Fig. 8 we note that SCIAMACHY nadir profile data does not pick up this anomalous behaviour of 2010 (also in red circles), the values in average are much lower than those in Fig. 9 and at most comparable to the rest of the years in Fig. 8 with variations in the range of 170 -100 DU.This could suggest that SCIAMACHY nadir profiles using plain L1 v8 data on its own may not be accurate enough (owing to remaining biases in L1 data, other instrumental issues like the instrumental coverage) to investigate inter-annual variability in Antarctic ozone depletion.Note that there is no SCIAMACHY data in early August south of latitude 70 • due to the low Sun.Furthermore note that the MSR v2 does a time dependent bias correction before assimilation takes place.However, in the region of 55 dataset (van der A et al., 2015) shown in Fig. 9 there are no significant deviations of ozone columns for 2003 (left panel).In a study by (Tilstra et al., 2012), it was shown from the retrieval of Absorbing Aerosol Index (AAI) using SCIAMACHY nadir spectra from 340-380 nm that strong jumps in the neighbouring AAI were observed from day to day in the years 2003, 2004 and 2008.These are exactly correlated with the days where the instrument was heated up to get rid of the ice layer affecting the infrared wavelengths and their degradation.Probably this operation on SCIAMACHY instrument affects all the versions of L1 data.We suggest that these instrumental throughput changes are the cause of the deviating ozone profiles for the years 2003-2004.
In Sect.3, it was shown that spectral corrections like shift and squeeze at the UV wavelengths can further improve the solar spectra by ∼ 4% (see Fig. 2).However improvement from this is much less than the expected degradation and other remaining potential biases for the L1 data which are increasing with the mission time.A preliminary analysis of this is shown in Fig. 11, where the mean ratio of observed to simulated reflectance spectra (R meas /R sim ) is plotted for the day of June 24 for years 2003-2011 in fading black to white colour.There is a strong deviation of this ratio at 265-290 nm and this becomes worse for later years, 2008-2011.A detailed study of this bias is beyond the scope of this paper.However, a L1 reflectance bias correction can significantly improve the quality of L1 data and will influence the validation and other results in this paper from their ozone retrievals.Furthermore applying these bias corrections would also make it meaningful to apply the spectral slit function corrections that can potentially improve the ozone retrievals further.Thus a detailed study of the effect of such bias corrections on the L1 data can be investigated against the quality of the ozone retrieval algorithm to better understand the quality of the SCIAMACHY nadir data.Appendix A: Averaging Kernels from OPERA retrievals using SCIAMACHY data The averaging kernel (AK) of a retrieval represents the measurement sensitivity with respect to the true state of the atmosphere.
The rows of the ozone profile averaging kernel matrix give the smoothing of the true profile as a function of the ozone retrieval layers.For an ideal retrieval, the curve of each row will peak at the nominal layer height with a spread that gives the vertical resolution of the retrieval.
In Fig. A1 we show an example of the AK for an individual OPERA ozone profile retrieval for a pixel on 2004/01/07.In
Figure 1 .
Figure 1.An example of the spectra of Earth radiance (solid blue line) and solar irradiance (dotted blue line) with the corresponding radiometric scale [photons/cm2/s/sr/nm] indicated on the left side of the figure, and reflectance (solid black line) with the corresponding scale [unit-less] indicated on the right side of the figure.The vertical dashed line in green separates the two channels and vertical grey lines
Figure 2 .
Figure 2. Relative residuals of daily SCIAMACHY L1 v8 solar spectra for Channels 1 and 2. G, O, D, S are Gain, Offset, Displacement/Shift and Squeeze respectively, where M is SCIAMACHY Measurement and S is Simulation (or modified model reference, see text).For visibility the residuals for ∼ 300 spectra are shown spread throughout the mission length from August 2002 to April 2012.
Figure 3 .
Figure 3. Temporal evolution of the slit function parameters shown for Channels 1 (top row) and 2 (bottom row).This is shown for the entire mission for 3466 days where Year 0 means 2002 and Year 8 means 2010.
2 2 .
At each given solar spectrum wavelength, λ i , this functional shape is numerically computed between [-1,1] nm centred at λ i .This shape can be manipulated with the SF parameters: additive constant (Offset, O), multiplication factor (Gain, G), Displacement of the peak along wavelength (Shift, D) and expansion and contraction of the spectral peak (Squeeze, S).The high resolution solar spectrum (Simulation) from Dobber et al. (2008) is modified with these four parameters to best match the SCIAMACHY measured solar intensity (Measurement, E(λ)).First spectral parameters, Atmos.Meas.Tech.Discuss., https://doi.org/10.5194/amt-2017-136Manuscript under review for journal Atmos.Meas.Tech.Discussion started: 9 June 2017 c Author(s) 2017.CC BY 3.0 License.
Figure 3
Figure 3 shows the temporal dependence of all the slit function parameters for Channel 1 (top row) and Channel 2 (bottom row) as density plots where the value for each wavelength and each day of the year is shown.The ranges of the values are to the right to each figure.The seasonal dependence of solar radiation is observed as expected in the parameter Gain (first column of the figure) for both channels.The other parameters do not show any seasonal dependence over the mission time.
Figure 4 .
Figure 4. Top-left: Measured reflectance spectra for the year 2003.The spectra are an average over all the pixels in the range of the geolocations around the sonde stations within a narrow tropic latitude bands from 10 • N to 10 • S. Top-right: Corresponding retrieved ozone profiles (in Dobson Units per layer) with different colours representing various versions of L1 SCIAMACHY data used (described in Section 2).The versions 7 with and without the degradation corrections almost overlap (blue and green lines) whereas version 8 with improved degradation correction has a visible difference in the ozone profile.Bottom rows are for the year 2009.Note the remarkable difference between v7 and later versions with degradation corrections.Here the later two versions almost overlap compared to the case where no degradation is taken into account showing the significant difference in L1 data between the three versions when degradation of the instrument is not taken into account.The number of pixels in each data set with corresponding uncertainties are listed in Table3.
Figure 5 .
Figure 5. Collocated geolocations indicating the location of the ozone sondes used in the validation for the dataset of L1 v8.Left: collocated geolocations for years 2003-2006 as labelled.Right: the same for years 2007-2011 as labelled.
Figure
FigureA2we show the validation results for the case where the sondes were convolved with the averaging kernel and the case where the sondes were not convolved.
median deviations for 2003 are smaller for v8 for Tr and NH compared to the older dataset versions, whereas for the SH the three different datasets give comparable deviations.The deviations for 2009 in v8 in the stratosphere are smaller for SH and NH than in the older datasets; the deviations are comparable in the Tr zone between all datasets.Comparison of the deviations between 2003 and 2009 show larger values for 2009 suggesting that the quality of L1 v8 data has degraded and is much worse than for earlier years (from comparison with 2003 v8 data).5.2 Validation of v8 for years 2003-2011 In Fig. 7 we show validation results for the entire dataset of v8 from the years 2003-2011 organized according to latitude bands with top panels showing results for SH, middle panels for Tr and bottom panels for NH.The left column show results for (early) years 2003-2006 and right column for (later) years 2007-2011.The solid lines correspond to the difference between satellite Atmos.Meas.Tech.Discuss., https://doi.org/10.5194/amt-2017-136Manuscript under review for journal Atmos.Meas.Tech.Discussion started: 9 June 2017 c Author(s) 2017.CC BY 3.0 License. the last two columns.From Fig. 7 we note that validation for all latitude bands show smaller deviations from the centre (zero) line for earlier years (left column), and for later years (right column) the agreement with sondes become worse especially for tropics and NH (right middle and bottom panels).The median deviations in the sixth and seventh columns of Table 4.3).It should be noted however that the large deviations in the tropospheric ozone [1000-100] hPa are systematically higher (than those for the for stratosphere) for all zones and the years even for v8, where such deviations are not observed for instance with GOME-2 nadir profiles(van Peet et al., private communication).This suggests that the quality of nadir SCIAMACHY L1 data is still poor and can be improved upon.It affects the lowest troposphere in the beginning of the mission and gets worse due to instrument degradation.This is still uncorrected in the UV wavelength range.Also observe the higher deviations in the year 2003 in SH and Tr (left upper and middle panels) compared to the rest of the years.This unique behaviour of 2003 is discussed in Sect.7.The deviations in SCIAMACHY v8 validation results above for the stratosphere can be compared for instance with GOME-2 validation results invan Peet et al. (2014) where they have used 16-layer pressure grid layer for retrievals.Their validation in the troposphere also showed deviations ranging from a few percent to ∼ −30% for the Northern Hemisphere, whereas we find the deviations to range from ∼ −10% to −45%) for year 2008 (blue line in right panel Fig.7.In the Southern Hemisphere however, we find these deviations for year 2008 (top right panel Fig.7) to range from a few percent to ∼ 20% which is more comparable to the range of a few percent to ∼ −15% in vanPeet et al. (2014).
• for the years 2003-2011 only for the months of (beginning of) August to (end of) November.The retrievals are separated in three latitude bands: [45 • S:55 • S], [55 • S:70 • S] and [70 • S:90 • S] shown in the top, middle and bottom rows of Fig. 8 respectively.The colour represents which year is labelled, showing early years (2003-2006) and later years (2007-2011) in the left and right columns of the figure, respectively.Each circle in the figure is the minimum value of the retrieved column for each day.The size of circle represents the number of pixels averaged per day and the range (minimum, maximum) of the number of pixels for all years and all latitude bands are listed in Table Atmos.Meas.Tech.Discuss., https://doi.org/10.5194/amt-2017-136Manuscript under review for journal Atmos.Meas.Tech.Discussion started: 9 June 2017 c Author(s) 2017.CC BY 3.0 License.ozone columns obtained from Brewer and Dobson spectrophotometers in the WOUDC dataset.The SCIAMACHY L1 dataset
Fig. 9
Fig. 9 is in line with the minimum occurring between 15 September to 15 October for the years 2003-2011.The lowest value in the SCIAMACHY ozone in Fig 8 occurs between 1 September and 15 October for latitude bands of 55 • S:70 • S and 70 • S:90 • S in the middle and bottom panels.The ozone columns reach a plateau in Fig. 8 from 1 November for latitude band 55 • S:70 • S where as this feature is not visible for the band 70 • S:90 • S which is more consistent with the Fig. 9.The overall level of the ozone column is at ∼ 150 in September for 70 • S:90 • S for the years 2004-2011 which is also qualitatively consistent with the MSR dataset.
• S:70 • S (mid-right panel of Fig.8we observe that the minimum ozone column for 2010 does exhibit the anomalous behaviour where it is higher compared to the rest of the years.This latitude region includes 70 • S, which samples the vicinity of the outer edge of the ozone hole, this result does show that SCIAMACHY nadir profiles could be used to complement the inter-annual variability studies.The discrepancy between the minimum of the total ozone column in regions 55 • S:70 • S and 70 • S:90 • S can be further investigated by carrying out a bias study of the L1 data (see Sect. 7).In Fig.10, the ozone profiles are shown for various latitude bands as in Fig.8.The median of the profiles for each latitude band and for each year is plotted as labelled with same colour coding as in Fig.8.In general the maximum of the ozone for all years decreases with the southernmost latitude bands going from top to bottom of each column in Fig.8.The peak value of the ozone (maximum) decreases from ∼ 30 DU at latitude band of 45 • S:55 • S to ∼ 25 DU at the latitude band of 70 • S:90 • S for all the years.Furthermore, the curve tends to be bi-modal at the southernmost latitude band with a minimum value of ozone (a value of zero DU) at a pressure of 100 hPa.This is the expected height of ozone depletion.The number of pixels for each year and for each latitude band are listed in Table6along with the median uncertainties in the ozone profile along with its spread.Columns 5-7 in Table 6 list the median number of iterations to reach conversion and the cost function that measures Atmos.Meas.Tech.Discuss., https://doi.org/10.5194/amt-2017-136Manuscript under review for journal Atmos.Meas.Tech.Discussion started: 9 June 2017 c Author(s) 2017.CC BY 3.0 License. the deviation of simulated and measured spectra at n-th iteration.What can be seen is the increasing number of iterations and corresponding increase in the cost function (thus worsening retrievals) for the southernmost latitude bands.7 Discussion In the previous section the ozone hole in the Antarctic was analysed using SCIAMACHY data for years 2003-2011 using geolocations south of 45 • .It is evident from Fig. 8 that the integrated ozone for the year 2003 (shown in black open circles) in the left panels are outliers.The retrieved columns deviate significantly in the latitude band 55 • S:70 • S (middle row) in the month of August and the deviations are significant for the same year in latitude band 70 • S:90 • S (bottom row) for all the months.Furthermore, the total ozone columns of the year 2004 are similar to those of 2003 compared to the other years in the bottom row of Fig. 8.This behaviour is reinstated further in Fig. 10, in the left -middle and bottom rows where the ozone profiles in the those bands for years 2003 and 2004 show a significant deviation from the rest of the years.In fact, the southernmost ozone profiles also show that the deviations are strong for the years 2004, 2005.From the assimilated and calibrated MSR v2 Atmos.Meas.Tech.Discuss., https://doi.org/10.5194/amt-2017-136Manuscript under review for journal Atmos.Meas.Tech.Discussion started: 9 June 2017 c Author(s) 2017.CC BY 3.0 License.We have performed nadir ozone profile retrieval using the OPERA algorithm for the latest complete SCIAMACHY dataset (v8) for almost the entire mission length from 2003 to 2011.Differences between datasets with and without degradation corrections (m-factors) were analysed in the wavelength range of 265-330 nm and show that the degradation correction including the scanangle dependence in the L1 v8 dataset gives the most smooth ozone profiles which is also reflected in their validation against ozone sondes.Retrieving the instrument slit function in the UV range for this L1 data set also gives improvement of a few percent in the solar data through the mission length.However, the measured reflectance spectra show that the degradation in v8 is still significant because the ratio of the measured to simulated reflectance spectra (calculated from the radiative transfer model used in OPERA) can range from ∼ 1.1 − 1.4 (seeFig.11, Sect.7).Furthermore the comparison between different L1 versions shows that v7 gives significantly worse ozone profiles especially later in the mission (2009) compared to v7 mfac and v8, where the profiles v7 show a double peak for the year 2009 and over-all reduced amount of ozone.Thus L1 v8 should be used for the nadir ozone profile applications of SCIAMACHY data.Using all L1 v8 data below 45 • S for years 2003-2011 we investigated the Antarctic ozone profile behaviour in the austral spring season.The daily minimum of the total ozone column from August to end of November shows a characteristic "V" shaped curve with the dip in mid-September to beginning of October at the southernmost latitude bands.This is consistent with other satellite data sets (for example,(van der A et al., 2015)).The outliers are the years[2003][2004], which can also be seen in the profiles for various latitude bands averaged over the whole year.The overall peak value of ozone reduces with southernmost latitude bands and a prominent minimum in the ozone profile with vanishing ozone concentration appears at a height of 100 hPa, as expected.
Fig. A2 ,
Fig. A2, we show the effect on the validation of including the satellite AK to the ozone sondes for the years 2003 (left panel) and 2009 (right panel).The validation results are clearly less noisy and smoother for the case where the AK was applied to the ozone sondes.
Figure 8 .Figure 9 .
Figure 8. Vertically integrated ozone columns (in DU) for the years 2003-2006 in left column for latitude bands 45 • S:55 • S (top row), 55 • S:70 • S (middle row), 70 • S:90 • S (bottom row).Each circle is the minimum value of daily retrieved quantity.The same is shown for years 2007-2011 in the right column.The time series are shown for months from 1 August -1 Dec as labelled in x-axis where a is August,s is September, o is October, n is November and d is December.Each dot is a median of the total column amount for that day corresponding to the latitude bands and its size represents the number of pixels (states) used in computing the median.
Table 1 .
SCIAMACHY Level 1 (L1) data characteristics Wavelength range [nm] Cluster number/Channel Integration time [ s] spectral resolution [nm] Integration time ranges from 0.125 s to 1 s depending on the which pixel it is. | 10,283 | 2017-06-09T00:00:00.000 | [
"Environmental Science",
"Physics"
] |
Novel Approach for Detecting Respiratory Syncytial Virus in Pediatric Patients Using Machine Learning Models Based on Patient-Reported Symptoms: Model Development and Validation Study
Background Respiratory syncytial virus (RSV) affects children, causing serious infections, particularly in high-risk groups. Given the seasonality of RSV and the importance of rapid isolation of infected individuals, there is an urgent need for more efficient diagnostic methods to expedite this process. Objective This study aimed to investigate the performance of a machine learning model that leverages the temporal diversity of symptom onset for detecting RSV infections and elucidate its discriminatory ability. Methods The study was conducted in pediatric and emergency outpatient settings in Japan. We developed a detection model that remotely confirms RSV infection based on patient-reported symptom information obtained using a structured electronic template incorporating the differential points of skilled pediatricians. An extreme gradient boosting–based machine learning model was developed using the data of 4174 patients aged ≤24 months who underwent RSV rapid antigen testing. These patients visited either the pediatric or emergency department of Yokohama City Municipal Hospital between January 1, 2009, and December 31, 2015. The primary outcome was the diagnostic accuracy of the machine learning model for RSV infection, as determined by rapid antigen testing, measured using the area under the receiver operating characteristic curve. The clinical efficacy was evaluated by calculating the discriminative performance based on the number of days elapsed since the onset of the first symptom and exclusion rates based on thresholds of reasonable sensitivity and specificity. Results Our model demonstrated an area under the receiver operating characteristic curve of 0.811 (95% CI 0.784-0.833) with good calibration and 0.746 (95% CI 0.694-0.794) for patients within 3 days of onset. It accurately captured the temporal evolution of symptoms; based on adjusted thresholds equivalent to those of a rapid antigen test, our model predicted that 6.9% (95% CI 5.4%-8.5%) of patients in the entire cohort would be positive and 68.7% (95% CI 65.4%-71.9%) would be negative. Our model could eliminate the need for additional testing in approximately three-quarters of all patients. Conclusions Our model may facilitate the immediate detection of RSV infection in outpatient settings and, potentially, in home environments. This approach could streamline the diagnostic process, reduce discomfort caused by invasive tests in children, and allow rapid implementation of appropriate treatments and isolation at home. The findings underscore the potential of machine learning in augmenting clinical decision-making in the early detection of RSV infection.
Introduction
Every winter, respiratory syncytial virus (RSV) causes acute lower respiratory tract infections in approximately 33.8 million children younger than 5 years worldwide [1].Approximately all children are infected at least once, and half are infected twice or more by the age of 24 months [2].Newborns and children with underlying medical conditions are particularly susceptible to severe infection [3][4][5].Therefore, reducing the number of RSV-infected patients is paramount for reducing the number of associated deaths.Consequently, there is an urgent need to develop a quick and accurate detection system for RSV infection [6].
Laboratory testing of RSV, including rapid antigen testing and polymerase chain reaction tests, provides reasonably accurate infection-related information [7].However, the collection of nasopharyngeal secretions, a necessary step for these tests, can cause discomfort in children.Furthermore, given that RSV test results rarely influence treatment decisions, these tests are not routinely conducted [8,9].On a positive note, RSV infections are not serious in most cases, and there is no active treatment, indicating that most patients can be treated at home under the constant supervision of parents or caregivers while taking precautions [10,11].Thus, remote identification of RSV infection allows patients to be cared for at home, preventing infection spread [12].This approach could also ease the burden on health care workers during epidemics by providing remotely procured information.Nonetheless, a home-based detection method that matches the accuracy of a rapid antigen test has yet to be recognized.
The symptoms and signs of RSV infection may help establish remote detection strategies.Studies that focused on the detection of RSV infection based on symptoms either lacked discriminatory accuracy or highlighted difficulties because of the diverse clinical manifestations of RSV infection [13][14][15][16][17][18].However, symptom onset of RSV infection may not appear as a cross-sectionally typical pattern at a specific time point but rather as a pattern that is diverse in characteristics, including the longitudinal aspects of symptoms.Particularly, the symptoms of RSV infection peak 4-5 days after infection and change with the increase or decrease in viral load [19,20].Dyspnea and other lower respiratory symptoms, including wheezing, moaning, and tachypnea, which are typical symptoms of severe RSV infection, occur when infected ciliated bronchial epithelial cells drop into the lower respiratory tract, thereby delaying the manifestation of upper respiratory symptoms [21][22][23].Contrarily, most studies linking RSV infection to overt symptoms used cross-sectional data based on symptoms at specific time points and did not consider the time course of symptom onset in the longitudinal profile of individual children.
Therefore, we propose that cross-sectional studies based on specific time points of signs and symptoms expressed with RSV infection appear to be unrelated to RSV infection, but this may not represent the unique disease trajectory of RSV infection in individual patients infected with RSV.Therefore, machine learning models based on symptom data structured to include longitudinal characteristics may enable highly accurate identification of viral infection.As the progression of symptoms can exhibit multiple patterns in each individual, considering aspects such as the size of patient bronchi and machine learning algorithms, which are already widely used to diagnose and classify diseases based on symptom characteristics, are equally suitable for identifying RSV infection [24][25][26][27].Here, we sought to leverage the longitudinal diversity of symptoms using machine learning based on patient-reported information, aiming to confirm the presence of RSV infection remotely at home with sensitivity and specificity comparable to those of rapid antigen testing [28].If this strategy is recognized, it will contribute to reducing the physical burden on children, saving medical costs, and preventing nosocomial infections.The purpose of this study is to develop and validate a machine learning-based RSV infection identification model using patient self-reported symptom information for outpatients.This study is also valuable as it is one of the few studies conducted in a cohort of outpatients with mild infections.
Overview
The results are presented according to the TRIPOD (Transparent Reporting of a Multivariable Prediction Model for Individual Prognosis or Diagnosis) statement [29].
Ethical Considerations
The study was conducted according to the Declaration of Helsinki and Japan's ethical guidelines.The institutional review board of Yokohama City Municipal Hospital approved the protocol (18-05-04), and we obtained informed consent from the patients' parents in the form of an opt-out clause.In this study, analysis was conducted using data that had been anonymized to ensure the privacy and confidentiality of participants.No compensation was provided to the participants.
Data Collection Setting
This observational retrospective cohort study involved patients aged ≤24 months who visited the pediatric or emergency department of Yokohama City Municipal Hospital between January 2009 and December 2015 and had RSV rapid antigen test results.According to the facility policy, all outpatients were required to fill an electronic template, and those who exhibited cold symptoms were subjected to RSV rapid antigen testing.These patients were extracted from a prospectively curated database and enrolled in this study.In this study, we used the immunochromatographic Quick-Navi RSV test (Denka Seiken Co Ltd) for nasopharyngeal swab fluid testing as the gold standard.Patients who presented weakly positive results in the rapid antigen test were excluded.
Data Preparation
An automated medical interview system with an electronic template was introduced to standardize the entry of clinical symptoms, and the patient's parents completed the form before the hospital visit.A group of highly trained general pediatric attending physicians with over 15 years of clinical experience created the template for entering symptoms and signs, allowing parents to select up to 3 symptoms per entry.Once a symptom XSL • FO RenderX was selected, the system presented additional questions based on the selection, with all responses except temperature being optional.Therefore, the status of each symptom, along with the number of elapsed days since symptom onset, was recorded as a categorical variable.Additional questions for each symptom were presented differently for each age group; all questions and options used are listed in Multimedia Appendix 1.The feature set was solely based on information from the electronic template and did not include any additional data from the examination or treatment.Based on medical insights, the symptoms introduced to the models were limited to cough, runny nose or nasal congestion, and wheezing, using only baseline characteristics and overall health status features.Statistical feature selection was not performed.
Experimental Process
We developed and evaluated a machine learning model that outputs binary information on RSV infection based on symptom and sign information.Users can obtain information about RSV infection status by entering symptoms to determine whether to seek medical attention.
Models
Random forest, extreme gradient boosting (XGBoost), and support vector machine models were used to determine appropriate machine learning algorithms.We used grid search in a hyperparameter space for all classifiers, optimizing the hyperparameters based on 10-fold cross-validation.The area under the receiver operating characteristic curve (AUC-ROC) was calculated for each model using the optimized hyperparameters.Finally, we selected the machine learning algorithm and its corresponding hyperparameters that performed best.
Model Performance Evaluation
Model performance was assessed through calibration and discrimination.Calibration was evaluated graphically using a calibration plot and Hosmer-Lemeshow test with 10 groups, where P<.05 indicated a poor model fit.The AUC-ROC, sensitivity, and specificity were used to evaluate model discrimination ability.Sensitivity and specificity were calculated using the Youden index, and performance was calculated under conditions, where 1 parameter was fixed to be equivalent to that of the rapid antigen test.Additionally, we calculated model discriminatory power based on the number of elapsed days since the onset of illness to assess the use of rapid patient isolation.The performance of a valid model should not be considerably worse even in short periods after disease onset.The discrimination metrics of the final model were evaluated to estimate stability using the 1000 times bootstrap method.This resampling technique, which involves generating multiple bootstrap samples and using them to train and test the model, offers a robust estimate of evaluation indices even with small samples by correcting for optimism in the model's performance [30].
To further assess the effect of additional symptom information on detection accuracy, we constructed a baseline model that classified patients based solely on the presence or absence of symptoms.As RSV infection is seasonal and its prevalence varies by season, a model excluding only the month of hospital visit was also created for comparison.
Interpretability Evaluation
We evaluated the interpretability of the final model using Shapley additive explanations (SHAP), calculated using an algorithm that mimics the Shapley value used in game theory to evaluate the relative importance of each feature on discrimination performance while considering interactions.Therefore, it was used to corroborate the presence of essential variables and interactions.R (version 4.2.0;R Foundation for Statistical Computing) was used for all analyses.
Patient Characteristics
Between January 2, 2009, and December 31, 2015, a total of 7362 patients underwent rapid antigen tests, and their parents provided the necessary information through an electronic template.Of these, 4182 patients who were aged 24 months or younger were included in the analysis.One patient with weak positive test results and 7 with inaccurate age information were excluded.Figure 1 depicts the process of patient exclusion, data selection, and missing value completion.Of the remaining 4174 patients, 619 (14.8%) were positive and 3555 (85.2%) were negative for RSV infection.Table 1 presents the demographic and clinical characteristics of patients.The details of the features used in the model are provided in Multimedia Appendix 1.
Calibration and Discrimination Ability of the Differential Models
The model with the XGBoost algorithm fit well visually, with good calibration in the Hosmer-Lemeshow goodness-of-fit test (P=.27).For discrimination ability, the AUC-ROC of the estimated model calculated with 1000 times bootstrap was 0.811 (95% CI 0.784-0.833).The sensitivity and specificity were 73.5% (95% CI 66.8%-79.4%)and 73.9% (95% CI 70.3%-77.2%),respectively.To validate the exclusion performance of the proposed model, the threshold was adjusted according to the discrimination performance of the rapid antigen test [28].When the sensitivity was set to 71.6% (95% CI 64.8%-77.8%),approximately 68.7% (95% CI 65.4%-71.9%) of the total patients were predicted to be negative for RSV infection.Patient samples were predicted to be positive (6.9%, 95% CI 5.4%-8.5%)when the specificity was set to 96.6% (95% CI 95.3%-97.7%).For the baseline model, which was considered to determine the effect of additional symptom information on performance improvement, the AUC-ROC was 0.766 (95% CI 0.739-0.792)for the model that excluded the month of visit (model 2), 0.703 (95% CI 0.677-0.731)for the model that considered only the presence of symptoms (model 3), and 0.521 (95% CI 0.484-0.556)for the null model that relied solely on age (model 4).The discrimination performance of the proposed model is summarized in Table 2.
Model prediction performance is calculated by adjusting thresholds under different conditions.The scores are based on 1000 times bootstrap, with means and 95% CIs of each sample.
Discrimination Based on Days Since Onset
Elapsed time was measured by counting the number of days since the initial onset of symptoms.A total of 46.8% (n=1952) of patients reported that symptoms such as cough, nasal discharge, or wheezing began within 3 days, and 77.4% (n=3231) of RSV-positive patients developed these symptoms within 8 days.The AUC-ROC was 0.721 (95% CI 0.628-0.815)for patients on the day of symptom onset and 0.746 (95% CI 0.694-0.794)and 0.779 (95% CI 0.749-0.808)for patients within 3 and 8 days of onset, respectively.To assess the robustness of these AUC-ROC values, 1000 bootstrap samples were used for computation, as depicted in the box and whisker plot in Figure 2.This figure illustrates the distribution of AUC-ROC values across the elapsed days since symptom onset, highlighting that the variance in AUC-ROC values diminishes as the sample size increases with more elapsed days.
Interpretability of the Final Model
Variables contributing to the prediction were examined using SHAP.The variable that primarily contributed to the relative prediction performance was the month of visit, followed by the number of days from the onset of cough and maximum body temperature.Based on the Beeswarm plot, the number of days from the onset of the cough variable exhibited bifurcation along the x-axis, indicating variability in its impact on the model's output.However, age or current body temperature did not show a clear trend between the feature value and SHAP (Figure 3).
Principal Results
We developed a machine learning tool that leverages longitudinal diversity of symptoms for the remote detection of RSV infection at home.The model that incorporated the time course of the onset of symptoms and their evolution (model 2) exhibited enhanced discriminative performance.However, our final model, including information on onset days and month of visit (model 1), performed effectively in terms of sensitivity and specificity.By adjusting the threshold based on 2 exclusion criteria set at the same level as the rapid antigen test, this model applied thresholds of 0.152 or 0.463, thereby performing an exclusion accuracy comparable to the standard RSV detection method and successfully identifying 75.6% of infected patients for exclusion.This approach could potentially reduce the need for further confirmatory testing-most positive cases comprised samples with values above the threshold based on high sensitivity.Samples with values below the threshold had almost no positives and could be considered negatives.Samples with values below the threshold based on high specificity comprised mostly of negatives, indicating that samples with values above the threshold could be considered as positives.Nevertheless, rapid testing should be continued for patients who fall within the defined thresholds.The consistency between the estimated positive probability and the actual percentage was confirmed using calibration.Hence, there is no reason to perform additional testing on samples labeled based on the 2 thresholds.
When considering the number of days since disease onset, our proposed model demonstrated an AUC-ROC of approximately 0.721 for patients who began to experience symptoms on day 1, which is the visiting day, showing particularly high accuracy when including patients with symptoms that emerged within 5 days.The results indicate that discriminatory ability improved within approximately 6 days; however, it was high even on the day of onset.
We examined the characteristics essential for model identification.The month of visit was the most crucial variable; as RSV infection has a clear seasonality, the model reflected a tendency for RSV infection to occur more frequently in January.Although the peak period of RSV infection currently tends to be earlier, the presence of seasonality itself cannot be negated [31].For the days of onset variable, there was a clear difference in the presence or absence of symptoms, but we could not identify a clear pattern in the number of days of symptom onset.This was also the case for age, suggesting that various interactions occurred among the characteristics.This observation suggests that the machine learning model ultimately captured additional symptom-related questions generated by skilled pediatricians and the complex onset patterns of patient underlying characteristics.This finding was also confirmed by the apparent difference in discrimination performance between a simple symptom-only model (model 3) and a model with additional symptom and patient information (model 2).
Comparison With Previous Studies
The model used here showed a higher discriminative performance than those used in other studies on symptom-based RSV infection detection.For children, a model with 80% sensitivity, 68% specificity, and an AUC-ROC of 0.66 has been reported [15].A model with 72.8% sensitivity and 73.2% specificity has also been reported; however, this model included x-ray and laboratory test results as features, which differed from the variables we used, where patients even outside the hospital could ascertain RSV infection by themselves [16].Therefore, the results are considered noteworthy in terms of identification accuracy.To the best of our knowledge, this is the first study to include outpatients and obtain symptom information through nonmedical personnel, which is in contrast to previous studies that involved inpatients or data obtained by health care professionals [13][14][15][16][17][18].
RenderX
Our results indicate that adequately accurate predictions can be acquired using machine learning and symptom information.The detection of RSV infection based solely on patient-reported symptoms is still in its early stages; however, our tool shows robust capabilities in distinguishing positive and negative results.The tool estimates RSV infection based on symptom data and progression entered by the patient's caregivers at home, indicating that the intervention of medical personnel is not required.This remote detection strategy could potentially reduce the risk of nosocomial infections and physical burden on children.Moreover, by applying this tool, isolation measures can be implemented before visiting a medical facility.This will allow RSV infections to be detected at home, reducing the need to visit hospitals, thereby preventing secondary and nosocomial infections and reducing the burden on health care providers, especially during an epidemic, and protect them from coinfection with RSV and severe acute respiratory syndrome coronavirus 2 [32].In addition, it may limit the spread of the virus in the community.
Another unique feature of this study is that the model was developed based on 3 primary symptoms: cough, nasal discharge, and wheezing, which were selected and extracted from the system based on the clinical manifestations of RSV infection.Therefore, although other symptoms were recorded in the electronic template, they were not used here based on clinical rationale.
Limitations
There are several limitations to this study.First, participant demographics were limited as this study was conducted at a single institution.However, it is worth noting that Japan's insurance system, which allows for free access, may mitigate potential economic-related biases in our findings.Additionally, the rapid antigen test may yield false positives when there is a low disease incidence [33,34].Thus, the final diagnostic result, used as the gold standard here, may vary.The RSV prevalence pattern may already be changing and combining this with a surveillance system would further improve accuracy.Furthermore, model calibration must be confirmed based on race to assess performance differences between races when introduced to populations with widely varying demographics.In this study, we used a single cohort for internal validation and did not perform external validation, necessitating further testing on additional data sets to confirm generalizability.
Future Research Directions
Future prospective studies are required to assess the generalizability of this algorithm to all patients because they may differ from the retrospective cohort used in this study.The model developed in this study is specific for RSV infection; however, similar methods may be used to construct models to detect infections with other viruses using their respective symptoms.
Conclusions
Our detection tool was based on patient-reported symptoms and basic attribute information; nevertheless, it effectively detected RSV infection.Furthermore, our findings highlight the necessity to develop machine learning models and support the use of structural data for capturing complex patterns for symptom-based detection of RSV infection.The presented model leverages the distinct temporal patterns of RSV symptoms, allowing accurate identification of the infection even at early stages and with symptom evolution.Health care providers can perform model analysis before an outpatient visit to direct infected patients to home treatment or an appropriate isolation cohort.Applying this model to other settings can validate a standardized and comprehensive approach to improve RSV infection detection at home, and it could then be applied to other viruses.
Figure 1 .
Figure 1.Sequence of steps in patient selection and data preprocessing.
Figure 2 .
Figure 2. Area under the receiver operating characteristic curve (AUC-ROC) for elapsed days in the proposed model.
Figure 3 .
Figure 3. Average Shapley additive explanations (SHAP) value for each feature from the top in order of importance.
Table 2 .
Discrimination performance under different conditions.
Specificity equivalent to that of the rapid antigen test (with threshold: 0.463)
a AUC-ROC: area under the receiver operating characteristic curve. | 4,982.6 | 2024-04-12T00:00:00.000 | [
"Medicine",
"Computer Science"
] |
LncRNA HIF1A-AS1 Regulates the Cellular Function of HUVECs by Globally Regulating mRNA and miRNA Expression
Background : Long non-coding RNA (lncRNA) hypoxia inducible factor 1 α -antisense RNA 1 ( HIF1A-AS1 ) serves critical roles in cardiovascular diseases (CVDs). Vascular endothelial cells (VECs) are vulnerable to stimuli. Our previous study revealed that knockdown of HIF1A-AS1 reduces palmitic acid-induced apoptosis and promotes the proliferation of human VECs (HUVECs); however, the underlying mechanism remains unclear. Material and Methods : Cell Counting Kit-8, flow cytometry, transwell invasion, and wound healing were applied to detect the function of HUVECs. Moreover, miRNA sequencing (miRNA-seq) and RNA sequencing (RNA-seq) were conducted to uncover its underlying mechanism. Quantitative Polymerase Chain Reaction (qPCR) was implemented to assess the accuracy of miRNA-seq. A co-expression network was generated to determine the relationship between differentially expressed miRNAs (DEmiRNAs) and differentially expressed genes (DEGs). Results : Knockdown of HIF1A-AS1 promoted the proliferation, migration, and invasion but reduced the apoptosis of HUVECs, and the overexpression of this lncRNA had the opposite effect. Numerous DEmiR-NAs and DEGs were identified, which might contribute to this phenomenon. Multiple target genes of DEmiRNAs were associated with cell proliferation and apoptosis, and overlapped with DEGs identified from RNA-seq. Finally, the network manifested that lncRNA HIF1A-AS1 moderated the function of HUVECs by not only regulating the expression of some genes directly but also by influencing a few miRNAs to indirectly mediate the expression of mRNAs. Conclusions : The results suggested that HIF1A-AS1 might regulate HU-VEC function by not only regulating the expression of some genes directly but also by influencing some miRNAs to indirectly mediate the expression level of mRNA.
Introduction
Cardiovascular disease (CVD) encompasses a variety of conditions that affect the heart and blood vessels, such as cerebrovascular disease, peripheral artery, and irregular heartbeat [1].It is a serious threat to human health, characterized by a high prevalence, high disability rate, and high mortality.The number of people who die of CVD and cerebrovascular disease annually is up to 15 million, ranking first among all causes of death.Vascular endothelial cells (VECs), which are located in the innermost part of blood vessels, are vulnerable to stimuli.The apoptosis of VECs is closely correlated with numerous cardiovascular diseases (CVDs) such as arteriosclerosis, thrombus formation, and plaque erosion [2,3].Due to the complex pathogenesis and serious complication, an increasing number of studies has focused on the etiology of VEC injury.However, the underlying mechanisms remains unclear, hindering the prevention and treatment of related diseases.Hence, it is necessary to identify novel apoptosis-related therapeutic targets in VECs for CVD.
In recent years, non-coding RNAs (ncRNAs), mainly circular RNAs, microRNAs (miRNAs), and long ncRNAs (lncRNAs), have been widely investigated [4].LncRNA is widely expressed and plays an essential role in numerous life activities such as regulation of the cell cycle and cell differentiation.The abnormal expression or function of lncRNA is closely related to the occurrence of human diseases including cancer, immune responses, and other related diseases [5,6].The biogenesis of lncRNA is associated with its specific subcellular localization and function.LncRNA is a potential biomarker that can be applied to clinical targeting and has potential therapeutic effects [7].
A number of noncoding RNAs, including lncRNAs and miRNAs, play pivotal roles in the progression of vascular diseases [8][9][10][11][12].Some of them are emerging as diagnostic biomarkers or therapeutic targets due to their specific role in some CVDs [13][14][15][16].LncRNA hypoxia inducible factor 1α-antisense RNA 1 (HIF1A-AS1), as one of three types of HIF1A antisense RNA, is located on the antisense strand of HIF1A of human chromosome 14, and its length is 652 nucleotides (nt) [17].This lncRNA is widely distributed in the bone marrow, appendix, and gall bladder, among other organs.Furthermore, the expression of HIF1A-AS1 shares a heterogeneous spatial distri-bution across normal human tissues [18].HIF1A-AS1 is highly upregulated in human VECs (HUVECs) by acriflavine, a DNA topoisomerase inhibitor, which may also damage HUVECs [19].Recent studies have demonstrated that this lncRNA might regulate the apoptosis and proliferation of vascular smooth muscle cells (VSMCs) [20].In addition, HIF1A-AS1 has pro-apoptotic and pro-inflammatory roles in Coxsackievirus B3-induced myocarditis by targeting miR-138 [21].The interaction between apoptotic proteins and HIF1A-AS1 plays an important role in the proliferation and apoptosis of VSMCs cultured in vitro, which might be involved in the pathogenesis of the thoracoabdominal aortic aneurysm [22].HIF1A-AS1 from exosomes could function as potential biomarkers for atherosclerosis [22].Our previous studies indicated that the suppression of HIF1A-AS1 can promote the proliferation and reduce the apoptosis of HUVECs induced by palmitic acid (PA) treatment [23], suggesting that this lncRNA might play critical roles in regulating HUVECs.However, the potential functions and regulatory mechanisms of this lncRNA in CVD related to HUVECs have not been fully elucidated.This study identifies novel therapeutic targets related to VEC apoptosis and provides new ideas for the treatment of CVDs, which is of great significance for the prevention and treatment of CVD.
In the current study, we elucidated the role of HIF1A-AS1 in HUVECs and its underlying mechanisms.The results could provide insights into potential research directions for CVD treatment in the future.First, as previously reported [24], we simulated cardiovascular occlusion by treating HUVECs with PA. qPCR confirmed the successful transfection of plasmids containing HIF1A-AS1 or short hairpin RNA (shRNA).Using flow cytometry, Cell Counting Kit-8 (CCK-8), transwell and wound healing assays, we found that this lncRNA promoted apoptosis and reduced proliferation, migration, and invasion.Moreover, miRNA sequencing (miRNA-seq) results showed that HIF1A-AS1 globally mediated the expression of miRNAs.Bioinformatics analysis indicated that multiple target genes of differentially expressed miRNAs (DEmiRNAs) were involved in cell metabolism and apoptosis.Subsequently, RNA sequencing (RNA-seq) and bioinformatics analysis was applied to identify differentially expressed genes (DEGs).Interesting, quite a few DEGs overlapped with the target genes of DEmiRNAs.Finally, a co-expression network showed the strength of the correlation between expression levels of DEmiRNAs and some initial DEGs.These data suggest that HIF1A-AS1 regulates the function of HUVECs by not only directly regulating the expression of some genes but also by influencing some miRNAs to indirectly mediate the expression of mRNA.The experimental process and relevant mechanism are schematically illustrated in Scheme.
Scheme.Schematic of the experimental process and related mechanism.
Cell Culture
HUVECs were obtained from the Shanghai Cell Bank of the Chinese Academy of Sciences (Shanghai, China).The cells were taken from a male, and cells from passages 3 to 8 were used for the experiments.The cells were cultured in Dulbecco's modified Eagle medium (Catalog Number: 30030, Gibco, Waltham, MA, USA) containing 10% fetal bovine serum, 100 µg/mL streptomycin, and 100 U/mL penicillin.Then the cells were incubated at 37 °C in a standard atmosphere (Thermo Fisher Scientific, Waltham, MA, USA) with 5% CO 2 .
Plasmid Generation, Lentivirus Package, and Transfection
HIF1-AS1 was amplified and subcloned into the EcoRI and BamHI restriction sites of pLVX-Puro 1.0 empty plasmid (Thermo Fisher Scientific).About 2 µg of plasmids was mixed with the lentivirus packaging plasmids pHelper 1.0, pHelper 2.0, and Opti-MEM according to a previous standard protocol [24].Subsequently, the lentivirus was diluted in fresh medium and incubated for 24 h and the cells were washed.We silenced HIF1A-AS1 in HUVECs using lentivirus-mediated shRNAs.The shRNAs targeting HIF1A-AS1 were synthesized by Genepharm (Shanghai, China).Transfection of plasmids containing shRNAs or blank plasmids were conducted using Lipofectamine 3000 (Catalog Number: L3000001, Thermo Fisher Scientific, Shanghai, China) according to the manufacturer's protocol.The cloning primers and shRNAs are presented in Table 1.
Flow Cytometry
The apoptosis of HUVECs was analyzed by flow cytometry using an Annexin V-Conjugated FITC Apoptosis Detection Kit (Catalog Number: BMS500FI-20, BD Biosciences, Franklin Lakes, NJ, USA).The cells were divided into NC, PAT, OE, sham-OE, sh, and sham-sh groups.Briefly, prepared cells were harvested after cultivating for 48 h, washed twice with phosphate-buffered saline (PBS) and incubated in the dark with Annexin V-FITC and propidium iodide (PI) for 30 min.Subsequently, the stained cells were detected with the MoFlo XDP flow cytometer (Catalog Number: V145577, Beckman, Brea, CA, USA) and Cell Quest 3.3 software (BD Biosciences).
Transwell Invasion
To perform the invasion assay, chambers were assembled in 24-well plates with 8 µM pore transwell inserts (Catalog Number: 353504, BD Falcon, Franklin Lakes, NJ, USA) coated with 50 µL Matrigel (diluted 1:4 in serum-free media).Treated cells (1 × 10 5 ) were added to the medium of the upper chamber.The invasive cells at the bottom of the implant were fixed in 4% paraformaldehyde and stained with 0.1% crystal violet.We captured images with a stereo microscope (Leica-M165 C, Wetzlar, Germany).Cell were counted under the TS100 microscope (Catalog Number: Eclipse TS100, Nikon Instruments, Shanghai, China).
Wound Healing
HUVECs (5 × 10 5 cells) were treated with the indicated drugs PA in a 6-well plate for 48 h.The next day, a 10 µL pipette tip was used to draw a linear or circular scratch wound in the confluent monolayer of cells.The cells were washed three times with PBS to remove the cell debris, and then cultured in fresh serum-free media for 12 h in 37 °C, 5% CO 2 incubator.Images of the wound were captured at 0, 24, and 48 h at 40× magnification, and ImageJ software (Mac v2.3.0,LOCI, University of Wisconsin, Madison, WI, USA) was used to measure the extent of the wound size in three wells per group.
RNA Extraction, Small RNA-seq and mRNA-seq
Total RNAs were extracted from HUVECs with TRIzol reagent (Catalog Number: 15596026, Invitrogen, Thermo Fisher Scientific).The absorbance of purified RNA at 260 and 280 nm and the A260:A280 ratio were measured using the NanoDrop ND-1000 spectrophotometer (Bio-Rad, Guangzhou, China).Six groups (NC, PAT, sham-OE, OE, sham-sh, and sh) were generated and two biological replicates were made.
Total RNAs (3 µg) from every sample were used to construct small RNA cDNA libraries with the Balancer NGS Library Preparation kit (Catalog Number: K02420-S, Gnomegen, San Diego, CA, USA).According to the manufacturer's protocol, the whole library was submitted to 10% native polyacrylamide gel electrophoresis.Then the corresponding strip into which miRNA was inserted was cut and eluted.These small RNA libraries were processed with NextSeq X-10 (Illumina, San Diego, CA, USA).The libraries were generated using purified mRNAs with the TruSeq Stranded Total RNA LT Sample Prep Kit (Catalog Number: RS-200-0012, Illumina, San Diego, CA, USA), according to the manufacturer's protocol.AMPure XP(Catalog Number: A63880, Beckman, USA) beads were used to select cDNA with a length of 350 to 400 base pairs.To collect the RNA-seq data, the Nextseq X-10 system (Illumina, San Diego, CA, USA) was employed.RNA-seq data were mainly analyzed with Cytoscape (version 3.0.2,https://cytoscape.org/) and Hisat2 software (Hisat-2.1.0,The Johns Hopkins University, Baltimore, MD, USA).
Identification of Conserved and Novel miRNA
The FASTX-Toolkit (version 0.0.13,http://hannonlab.cshl.edu/fastx_toolkit/)was applied to process raw reads to obtain reliable clean reads.RNAs <18 or >35 nt in length were discarded from further analyses in view of the length of the mature miRNA and adapter lengths.Subsequently, the Rfam database (version 12.0, http://rfam.xfam.org/)was used to search the highquality clean reads.Hereafter, the kept unique sequences were aligned against the miRBase database [25] by using Bowtie, with one mismatch allowed.The aligned small RNA sequences were matched by conserved miRNAs, and the unmatched sequences might be potential candidates for new miRNAs.Finally, the unique sequences were mapped to the reference genome sequence (GRCH38) of humans by the miRDeep algorithm to identify novel miRNAs [26].
qPCR
To validate the miRNA-seq data, qPCR was conducted.B-cell lymphoma 2 (Bcl-2)-associated X protein (BAX ) and matrix metalloproteinase 1 (MMP1), key factors of HUVEC apoptosis and proliferation, were also validated by qPCR.The primers used are presented in Table 1.The PCR experiments were performed with the following conditions: pre-denaturation at 95 °C for 1 min, 40 cycles of denaturing at 95 °C for 15 s, annealing at 60 °C for 30 s, and elongation at 72 °C for 40 s.The results were calculated with the 2 −∆∆Ct method [27].
Western Blot Analysis
Western blotting (WB) was performed according to standard methods.
Total cell lysates were made in 1× sodium dodecyl sulfate buffer.
Equal protein concentrations were resolved by sodium dodecyl sulfate-polyacrylamide gel electrophoresis, and the proteins were electrotransferred to PVDF membranes (Bio-Rad).GAPDH (Sigma, St. Louis, MO, USA) was used as a loading control.Antibodies against BAX (1:2000) and MMP1
Bioinformatics Analysis
To detect the miRNA expression profile of the miR-NAs identified, the frequencies of miRNA counts were normalized to transcripts per million (TPM) with the following formula: normalized expression = actual read count/total read count × 10 6 .The strict standard of Padj <0.01 and log 2 (fold change) >1 or log 2 (fold change) <-1 indicated statistically significant DEmiRNAs.The expression profiles of mRNAs were normalized to fragments per kilobase of exon model per million mapped fragments.DEGs were obtained using the same method.
To predict the gene function of candidate genes or DEGs, Kyoto Encyclopedia of Genes and Genomes (KEGG) and Gene Ontology (GO) analyses were conducted using the DAVID bioinformatics database [28].Coexpression networks were generated by calculating the Pearson's correlation coefficient (PCC) for the expression levels of candidate genes or DEGs.To display the co-expression networks, Cytoscape (version 3.0.2) was employed.
Statistical Analyses
All values are presented from independent experiments done in triplicate as the mean ± standard deviation.For comparison, GraphPad Prism 7 (GraphPad Software Inc., San Diego, CA, USA) was used, and the significance of differences between the means was determined by the Student's t-test or one-way analysis of variance.p < 0.05 was considered statistically significant result.
LncRNA HIF1A-AS1 Regulates the Proliferation of HUVECs
We constructed a model of cardiovascular occlusion by treating HUVECs with PA [23].To investigate how HIF1A-AS1 regulates HUVEC function, we overexpressed or knocked down this lncRNA in PA-treated HUVECs (Fig. 1A), and analyzed its expression level by qPCR (Supplementary Figs.1,2).Compared with the NC group, the overexpression and knockdown of HIF1A-AS1 were successful.Then, we obtained six groups of cells with different treatment strategies, namely NC, PAT, PAT + sham-OE, PAT + OE, PAT + sham-sh, and PAT + sh groups.
The effect of HIF1A-AS1 on the proliferation of HU-VECs was determined with a CCK-8 kit.A significant time-dependent increase in proliferation was found among the six groups.CCK-8 assays showed that treatment of HU-VECs with PA for 72 h resulted in a 45.8% reduction in HUVEC survival.Silencing of HIF1A-AS1 significantly increased the absorbance values at all time points (24, 48, 72 h).In addition, OE of HIF1A-AS1 led to a 17.1% reduction of its basal inactivation (45.8%) at 72 h (Fig. 1B).The results of the CCK-8 assay indicated that silencing of HIF1A-AS1 promoted the proliferation of HUVECs, and OE led to the reverse effect.Flow cytometry was conducted to count the number of apoptotic cells.After treatment of HUVECs with PA, the number of apoptotic cells increased more than 3 times (Fig. 1C).Indeed, the apoptosis of HUVECs has been identified as a crucial factor in the pathogenesis of various CVD processes [29,30].Thus, preventing the apoptosis of HU-VECs may be a novel strategy for the treatment of CVD.Fig. 1C shows that silencing HIF1A-AS1 led to about 35.5% recovery, and OE led to 28.5% promotion of basal apoptosis induced by PA.These results suggest that the lncRNA HIF1A-AS1 plays a critical role in regulating the apoptosis and proliferation of HUVECs.
The expression of BAX and MMP1, key regulators of HUVEC apoptosis and proliferation, was analyzed by qPCR (Fig. 1E,F) and WB (Fig. 1G and Supplementary Figs.3,4,5).Quantitative analysis of the WB results was also performed (Supplementary Fig. 6), which was consistent with the corresponding qPCR results.Compared with the sham-OE group, the expression of BAX and MMP1 was significantly increased in the OE group.These data indicate that OE of this lncRNA promotes PA-induced apoptosis and reduces the proliferation of HUVECs.
LncRNA HIF1A-AS1 Regulates the Migration and Invasion of HUVECs
The transwell migration assay was performed to determine whether OE of HIF1A-AS1 can mediate the migration and invasion of HUVECs.In this assay, the cell invasion ability of PA-treated HUVECs decreased about 49.8% compared with blank groups.A 27.8% recovery was observed in invading cells upon silencing of HIF1A-AS1 compared to sham-sh cells, suggesting that silencing this lncRNA can enhance the migration ability of HUVECs.Conversely, OE of HIF1A-AS1 resulted in another 35.3% reduction, clearly indicating the weakened migratory ability of PA-treated HUVECs (Fig. 2A,B).
To further estimate the migration effect of HIF1A-AS1 lncRNA on HUVECs, we conducted a wound healing assay.As shown in Fig. 2C,D, the migration ability was weakened after OE of HIF1A-AS1 compared with the sham-OE group.By contrast, HUVECs migrated for a longer distance in the pore plate extracted from HIF1A-AS1 transgenic cells than from the sham-sh group.These results suggested that this lncRNA successfully inhibited HUVEC migration.
LncRNA HIF1A-AS1 Globally Regulates the Expression of miRNAs
To investigate how HIF1A-AS1 regulates the expression of miRNAs in HUVECs, 12 small RNA libraries (NC-1, NC-2, NC-3, PAT-1; PAT-2, PAT-3, PAT-sham-OE-1, PAT-sham-OE-2; PAT-sham-OE-3, PAT-OE-1, PAT-OE-2 and PAT-OE-3) were constructed for miRNA-seq.More than 141.2 million reads were generated, with approximately 11.8 million sequence reads per sample.Comparing all reads to the human reference genome, it was found that 91.2% of the reads were successfully mapped to the reference genome (Table 2).More than 852 conserved miR-NAs and about 29 novel miRNAs were identified from the boxplots of the 12 samples, no obvious differences were found in these groups (Fig. 3A and Supplementary Table 1).The above results confirmed the reliability and stability of miRNA-seq.
Identification of DEmiRNAs
MiRNAs act as key post-transcriptional regulators in multiple cellular biological processes such as proliferation, differentiation, apoptosis, invasion and migration [31].Interestingly, we found that OE of HIF1A-AS1 markedly regulated the expression of miRNAs.The results of the Pearson correlation data are shown in Fig. 3B.A heat map was generated to reflect the detailed alterations of miRNAs (Fig. 3C).In addition, the volcano plots in Fig. 3D compared these log 2 fold change values with their respective p (-log 10 ) values to obtain the distributions of both upregulated and downregulated DEmiRNAs.By setting a strict threshold, a total of 59 statistically significant DEmiRNAs were identified in the OE vs. sham-OE group (Fig. 3E and Supplementary Table 2).Furthermore, the Venn diagram showed 10 reliable core miRNAs across two groups between PAT vs. NC and OE vs. sham-OE (Fig. 3F).These findings indicate that HIF1A-AS1 is associated with the expression levels of some miRNAs in HUVECs.To confirm the accuracy and reliability of the miRNA assays, the expression levels of certain DEmiRNAs, including four upregulated miRNAs (hsa-miR-1298-5p, hsa-miR-30c-5p, hsa-miR-27b-5 and hsa-let-7a-5p) and four downregulated miRNAs (hsa-miR-4664-3p, hsa-miR-769-5p, hsa-miR-106b-5p and hsa-miR-548o-3p), were further validated by qPCR.All of the miRNAs were randomly selected from those 59 DEmiRNAs.The primers used are presented in Table 1.The results were in accordance with the miRNAseq results (Fig. 3G).
Target Genes of DEmiRNA Are Mainly Related to Apoptosis and Metabolism
MiRNAs commonly exert their functions through binding to complementary target sites from the target genes.
DEmiRNA was obtained from OE vs. sham-OE groups.The DEmiRNA target genes, also called candidate genes, were successively identified by miRBase.As shown in Supplementary Table 3, we obtained a total of 2030 and 2401 predicted targets of the upregulated and downregulated DEmiRNAs, respectively.For the 28 upregulated DEmiRNAs, hsa-miR-193b-3p was found to potentially target the most genes, with a number of 334.For the 31 downregulated DEmiRNAs, hsa-miR-5088-5p possessed the most targets, with a number of 1087.
It was obvious that OE of HIF1A-AS1 globally affected miRNA expression in HUVECs.Consequently, bioinformatics analysis was conducted to identify the key functions in which all candidate genes were involved.GO analysis showed that target genes were mainly enriched in the single-organism process (GO:0044699), cellular process (GO:0009987), and some metabolic processes, especially the organic substance metabolic process (GO:0071704).Furthermore, cellular component analysis indicated that the target genes were mainly enriched in the intracellular regions and shared a heterogeneous spatial distribution across the entire cell.In addition, molecular function analysis showed that most of them performed binding activities and enzymatic reactions (Fig. 3H).KEGG enrichment analysis of the target genes was performed to gain further insight into their functions (Supplementary Table 4).The results showed that most of the candidate genes were enriched in pathways associated with metabolism including metabolic pathways (ID: hsa01100), glycerophospholipid metabolism (ID: hsa00564), glycosaminoglycan degradation (ID: hsa00531), galactose metabolism (ID: hsa00052), hematopoietic cell lineage (ID: hsa04640), and glycine, serine, and threonine metabolism (ID: hsa00260) (Fig. 3I).It can be proved from the previous discussion that HIF1A-AS1 overexpression in HUVECs can produce numerous complex changes, indicating that it plays an important role in CVD.
mRNA Expression Profiles and Their Bioinformatics Analysis
RNA-seq experiments were done to further explore the molecular mechanisms by which HIF1A-AS1 regulates HUVECs.More than 1.20 billion pairend reads were generated, corresponding with ~100 million sequence reads per sample.Using Hisat2 software, >85.9% of clean reads were successfully mapped against the current human reference genome (GRCH38).The ratio of multiple mapped reads was less than 6.1% (Table 3).There was no significant difference among these groups in the boxplots (Fig. 4A).
Subsequently, the RNA-seq results revealed that OE of HIF1A-AS1 regulated the expression for both miRNA and mRNA.Pearson's correlation data are shown in Fig. 4B.The heat map reflected the DEGs (Fig. 4C).By adopting the same method, a total of 196 DEGs were obtained from the sham-OE vs. OE groups (Supplementary Table 5), with 122 upregulated and 74 downregulated, respectively (Fig. 4D,E).We observed six reliable core mRNAs across two groups (PAT vs. NC and OE vs. sham-OE) from the Venn diagram (Fig. 4F).qPCR was conducted to verify the four upregulated mRNAs (TUBB3, ANGPTL4, ISG15 and IFI6) and four downregulated mRNAs (HIST4H4, HIF1A, HMGCS1 and BBX ) (Fig. 4G).The qPCR results were in accordance with that of mRNA-seq.Interestingly, 23.5% DEGs overlapped with DEmiRNA target genes.The overlapping DEGs are listed in Supplementary Table 6.According to previous reports [27][28][29], both mRNA and miRNAs had the ability to regulate the apoptosis and proliferation of HUVECs.To further explore the role of miRNA in HUVECs, we conducted this study.
Since some DEGs overlapped with DEmiRNA target genes, bioinformatics analysis was conducted to identify the biological function of DEGs.GO enrichment analysis showed that most of the enriched pathways were associated with multiple metabolic process, signaling pathways, and apoptosis with high confidence.Remarkably, the outcome also highlighted large-scale alterations in a variety of metabolic pathways and apoptosis processes upon the elevation of HIF1A-AS1 levels.Similarly, the results also demonstrated that a larger number of DEGs were related to extracellular matrices with overlapping distributions, and others were restricted to particular cellular loci.The majority had binding functions and enzymatic activity, especially transferase activity, with a fraction of them having other housekeeping functions (Fig. 4H and Supplementary Table 7).
The most significantly enriched KEGG pathways are shown in Fig. 4I and Supplementary Table 8.The results showed that fructose and mannose metabolism (ID: hsa00051), synthesis and degradation of ketone bodies (ID: hsa00072) and butanoate metabolism (ID: hsa00650) were the most significant pathways for enrichment.KEGG results showed that some DEGs were more uniformly enriched in tumor-related pathways such as p53 signaling pathway (ID: hsa04115), renal cell carcinoma (ID: hsa05211), Vitamin B6 metabolism (ID: hsa00750), pathways in cancer (ID: hsa05200), malaria (ID: hsa05144), and caffeine metabolism (ID: hsa00232).The main signaling and metabolism pathways determined by KEGG analysis will provide further insight into future research directions on mRNA.
Integrative Analysis of DEmiRNA and mRNA Expression
Generally, miRNAs have the capacity to recognize and bind to complementary 3'-untranslated regions of target mRNAs, which can lead to the degradation or transcriptional repression of mRNAs [32].To explore the relationship between DEmiRNA and DEGs, a co-expression network was generated by calculating the PCC for the expression levels of DEGs and DEmiRNAs.It showed a close correlation between the expression levels of DEmiRNAs and DEGs, which were enriched in some initial pathways (Fig. 5).A total of 37 DEmiRNAs and 33 DEGs were filtered into the co-expression network complex.The network manifested that lncRNA HIF1A-AS1 mediated the function of HUVECs by not only regulating the expression of some genes directly but also influencing a few miRNAs to indirectly mediate the expression level of mRNA.These findings may explain the underlying mechanism of HIF1A-AS1 in CVD.
Discussion
Most CVDs are related to the apoptosis of VECs, which is the main form of vascular injury [33].Previous studies have shown that the broken balance between VECs apoptosis and proliferation markedly contributes to the pathogenesis of CVD [23,34].An increasing number of studies has shown the critical effect of lncRNAs on regulating the proliferation and apoptosis of VECs in CVD [3,35,36].
Only about 2% of sequences in the human genome possess the ability of encoding proteins.Accumulating evidence has revealed that lncRNAs are related to human diseases as a biomarker or therapeutic target [37,38].Postnatally, lncRNAs have attracted a lot of attention due to their variety of biological roles including cell cycle control, cell proliferation, apoptosis, transwell invasion, embryonic development, and carcinogenesis by mediating the gene expression at the transcriptional, splicing, transportation, and translational levels [39,40].
MiRNAs involved in mRNA degradation or translation inhibition [41] are a group of evolutionarily conserved ncRNAs about 20-22 nt in length from hairpin pre-miRNA precursors [42].Accumulative evidence has revealed that miRNAs, alone or in combination with lncRNAs, are involved in regulating specific gene expression at the translation or transcription level.Then they can alter cell signaling pathways associated with different physiological and pathological processes [43][44][45][46].
A study found that HIF1A-AS1 TFR2 forms triplexes with EPH receptor A2 (EPHA2) and adriamycin (ADM) double-stranded DNA under regular and triplex-stabilized conditions upon DNA hairpin formation.Increasing the expression of HIF1A-AS1 can inhibit the expression of EPHA2 and ADM, whereas the downregulation of HIF1A-AS1 produces opposite results.These results suggest that the trimer formation region can mediate EPHA2 and ADM inhibition [47].Another study showed that the HIF1A-AS1 was significantly increased in gemcitabine (GEM)-resistant pancreatic cancer cells.HIF1A-AS1 enhanced the GEM resistance of pancreatic cancer cells by upregulating the expression of HIF1α and promoting glycolysis.HIF1A-AS1 may be a new therapeutic target for GEM resistance of pancreatic cancer in the future [48].In our study, we overexpressed or knocked down this lncRNA in PA-treated HU-VECs to explore how HIF1A-AS1 efficiently regulates the function of HUVECs.
Another study, which was published recently, showed that some lncRNAs regulate various cellular processes by acting as competing endogenous RNAs (ceRNAs) and binding proteins.For example, HIF1A-AS1, acting as a ceRNA, absorbed miR-204 to evaluate Suppressor of Cytokine Signaling 2 expression in cardiac function [21].This lncRNA participates in the regulation of proliferation, apoptosis, and the activity of the extracellular matrix proteins of VSMCs [18,49].Abundant evidence has indicated that this lncRNA might participate in the development of CVD by regulating the PA-induced apoptosis of HUVECs [23].However, the molecular mechanism by which HIF1A-AS1 interacts with miRNAs and mRNA and their regulatory roles of pathogenesis are unclear.HIF1A-AS1 has potential as a novel therapeutic target in CVD, but underlying information about the regulatory mechanisms in HUVECs is lacking.
It is known that miRNAs are involved in the progression and pathogenesis of VECs [50].In this study, it was found that OE of HIF1A-AS1 reduced the cellular growth rate and led to the robust apoptosis of HUVECs.We further studied the molecular mechanism of this phenomenon using miRNA-seq and RNA-seq.More than 852 conserved miRNAs were identified and about 29 novel miRNAs were found by miRNA-seq in the sham-OE group.When HIF1A-AS1 was overexpressed, the expression levels of some miR-NAs markedly changed, indicating that this lncRNA may play a critical role in miRNA-based therapies.The target genes of those DEmiRNAs were successively predicted by miRase.Additionally, multiple target genes of the DEmiR-NAs were associated with the apoptosis, proliferation, and migration of HUVECs, suggesting that OE of HIF1A-AS1 could inhibit proliferation and promote the apoptosis of HUVECs by mediating miRNA expression.As previously reported, lipids have important functions in maintaining normal physiological cellular functions [51].Glycosaminoglycan can promote wound healing.The administration of d-galactose to animals decreases the proliferation of cells and reduces the migration and survival of new neurons in the granule cell layer [52].Researchers have found a potential involvement of the glycine-serine-threonine metabolic axis in longevity and related molecular mechanisms [53].Thus, DEmiRNA may serve a regulatory role in the molecular functional analysis of HUVECs.Subsequently, RNA-seq was performed to identify the DEGs.Cluster of differentiation, which leads to endothelium apoptosis, was not identified.Individual differences may be responsible for this unusual phenomenon.Many DEGs were found to overlap with the target genes of DEmiRNA.Furthermore, both mRNA and miRNAs could regulate the apoptosis and proliferation of HUVECs [54][55][56].Therefore, HIF1A-AS1 has the ability to regulate the expression of some miRNAs, which could target some apoptosis-related genes by degrading mRNAs or inhibiting their translation.The crosstalk among miRNAs, lncRNAs, and mRNA shows a complex network of gene expression regulation [57].Hence, in the present study, a co-expression network was systematically constructed to explore the relationship among lncRNA HIF1A-AS1, DEmiRNA, and DEGs.The findings revealed that the expression levels of the DEmiRNAs was tightly linked to the apoptosis-related DEGs.However, this network has not been systematically validated, which limits the comprehensive understanding of the mechanisms underlying the role of HIF1a-AS1 in the treatment of CVD.In addition, accumulative evidence has indicated that ANGPTL4 is directly correlated with the risk of CVD, especially atherosclerosis [58].SERPINE1 may serve as a potential therapeutic target or new biomarkers in acute myocardial infarction [59].Interferon Alpha Inducible Protein 6, which is a mitochondrial localized antiapoptotic protein, contributes to promoting the metastatic potential of certain cancer cells through mitochondrial reactive oxygen species [60].However, in this co-expression network, the expression levels of these CVD-related genes were tightly related to certain miRNAs.Therefore, HIF1A-AS1 can modulate the expression of DEGs by mediating miRNA expression.The present study reveals a novel mechanism by which HIF1A-AS1 regulates the apoptosis of HUVECs.
Conclusions
In summary, our study showed that HIF1A-AS1 regulated HUVEC function by not only regulating the expression of some genes directly but also influencing some miR-NAs to indirectly mediate the expression level of mRNA, indicating that it may play a key role in the pathogenesis and progression of CVD.The current study also provides some new insights and directions for the prevention and treatment of CVD.Although the clinical applications need to be further explored, these results additionally provide insight into the molecular mechanisms by which HIF1A-AS1 affects HUVECs and a scientific experimental basis for treating CVD.Thus, it is feasible that the co-expression network could be applied for the prevention, diagnosis, treatment, and prognosis of CVD.However, further studies are being conducted to more systematically elucidate the role of HIF1A-AS1 in CVD and further determine the potential clinical role of the co-expression network.
Fig. 1 .
Fig. 1.Vector construction and HIF1A-AS1 affects proliferation and apoptosis in vitro.(A) qPCR was used to verify the successful construction of the vector.(B) The CCK-8 assay was conducted to evaluate the cell proliferation of six treatment groups at 0, 24, and 48 h.The data are presented as the percentage relative to control cells and presented as the mean ± standard deviation (SD) of three replicates.(C) Apoptosis was detected using Annexin V-fluorescein isothiocyanate staining coupled with flow cytometry.Every group had three parallel controls.The upper left, upper right, and lower right quadrants represent necrotic, late apoptotic, and early apoptotic events, respectively.(D) Total percentage of apoptotic HUVECs in each treatment group were quantified with the data presented as the mean ± SD of three independent experiments.(E) qPCR was used to analyze expression of the pro-apoptotic protein BAX.(F) qPCR was conducted to detect the expression of migration-related protein MMP1.(G) WB was performed to assess the expression of BAX and MMP1.GAPDH was used as a loading control for WB.Statistical analysis was carried out using one-way ANOVA followed by Tukey's post hoc test.*p < 0.05, **p < 0.005, ***p < 0.001, ****p < 0.0001.
Fig. 2 .
Fig. 2. HIF1A-AS1 affects migration and invasion in vitro.(A) & (B) The wound healing assay was used to detect the relative cell migration in six groups of cells, scale bar = 200 µM.Quantitative analysis of wound healing was performed for three fields.The migration capability of HIF1A-AS1 OE cells was significantly decreased.By contrast, HIF1A-AS1 shRNA markedly enhanced cell migration compared to the sham-sh group.(C) & (D) The wound healing assay showed that HIF1A-AS1 clearly decreased the invasion of the cells, whereas HIF1A-AS1 silencing showed the opposite effect.The scale bar = 100 µM.The quantitative data of the transwell assay were obtained from five fields.Values shown are the mean ± SD from three independent experiments, **p < 0.005, ***p < 0.001.
Fig. 3 .
Fig. 3. Exploration of DEmiRNAs and functional analysis.(A) The boxplots of the 12 samples miRNAs, (s) stands for sample, n = 12.(B) OE of HIF1A-AS1 could markedly regulate the expression of miRNAs.The results of Pearson's correlation data were presented.(C,D) Heat map (C) and Volcano plot (D) of DEmiRNAs expression profiles between OE and sham-OE.(E) The number of upregulated and downregulated DEmiRNAs among NC vs. PAT, OE vs. sham-OE, and OE vs. PAT.(F) Venn diagrams of the DEmiRNAs identified in different comparisons.Data are presented as the mean ± SD. (G) qPCR validation of certain DEmiRNAs identified by miRNAsequencing in the OE and sham-OE groups.Statistical analysis was conducted by the Student's t-test, and data are presented as the mean ± SD and of experiments conducted in triplicate.(H,I) GO (H) and KEGG (I) pathway enrichment analyses of target genes from sham-OE vs. OE.DEmiRNAs, differentially expressed miRNAs; NC, normal control group; PAT, PA-treated HUVECs; PAT + sham-OE,
Fig. 4 .
Fig. 4. Exploration of DEGs and functional analysis.(A) Boxplots of the 12 samples RNA-seq; no significant difference was found among these groups.(B) Subsequently, the RNA-seq results revealed that OE of HIF1A-AS1 not only regulate the expression of miRNAs but also genes.(C,D) Heat map (C) and Volcano plot (D) of DEGs expression profiles between the OE and sham-OE groups.(E) The number of upregulated and down-regulated DEGs among NC vs. PAT, OE vs. sham-OE and OE vs. PAT groups.(F) Venn diagrams of the DEGs identified in different comparisons.Data are presented as the mean ± SD. (G) qPCR were carried out for validation of certain DEGs identified by mRNA-seq in the OE and sham-OE groups.Statistical analysis was conducted by the Student's t-test and data are presented as the mean ± SD and experiment was performed in triplicate.(H,I) GO (H) and KEGG (I) pathway enrichment analyses of DEGs from sham-OE vs. OE.DEGs, differently expressed genes.Up-regulated mRNA: TUBB3, ANGPTL4, ISG15, IFI6; down-regulated mRNA: HIST4H4, HIF1A, HMGCS1, BBX.*p < 0.05, **p < 0.01.
Fig. 5 .
Fig. 5. Interaction network analysis of DEmiRNAs and DEGs associated with some important pathways.Network analysis on the basis of PCCs for DEmiRNAs and DEGs enriched in 'negative regulation of endothelial cell apoptotic process' (GO:2000352), 'cellular carbohydrate metabolic process' (GO:0044262), 'type I interferon signaling pathway' (GO:0060337) and so on.Circular nodes represent DEmiRNAs and rectangular nodes signify DEGs, while these solid lines represent significant correlations between DEmiRNAs and DEGs.The red lines represent negative correlation, while the blue lines represent positive correlation.p < 0.01 and PPC ≥ 0.06 indicated a statistically significant correlation.PCC, Pearson's correlation coefficient; miRNA, microRNA. | 7,996.8 | 2022-12-21T00:00:00.000 | [
"Medicine",
"Biology"
] |
Regularization Total Least Squares and Randomized Algorithms
: In order to achieve an effective approximation solution for solving discrete ill-conditioned problems, Golub, Hansen, and O’Leary used Tikhonov regularization and the total least squares (TRTLS) method, where the bidiagonal technique is considered to deal with computational aspects. In this paper, the generalized singular value decomposition (GSVD) technique is used for computational aspects, and then Tikhonov regularized total least squares based on the generalized singular value decomposition (GTRTLS) algorithm is proposed, whose time complexity is better than TRTLS. For medium-and large-scale problems, the randomized GSVD method is adopted to establish the randomized GTRTLS (RGTRTLS) algorithm, which reduced the storage requirement, and accelerated the convergence speed of the GTRTLS algorithm.
Introduction
In practical problems, many discrete ill-conditioned problems arising from many different fields of physics and engineering can be reduced to solving linear equations in the form of Ax ≈ b.The methods used commonly are least squares (LS) [1] and total least squares (TLS) [2,3].However, these kinds of problems are often ill-conditioned, such as the first kind of integral equations [4,5].In order to reduce the serious instability caused by the problems themselves, regularization treatment [6][7][8][9][10][11] becomes an effective method, that is, replacing the original ill-conditioned problem with an adjoining well-conditioned one, whose solution is called a regularized solution to approximate the true one.We know that Tikhonov regularization is one of the common methods, which is widely used in the industrial field [6].For example, Tikhonov regularized TLS (TRTLS) proposed by Golub, Hansen, and O'Leary can be used to approach the true solution.During the process, the bidiagonalization technique is used.It is shown that the ideal approximation solution cannot be obtained by the truncation singular value method in some practical problems.The total least squares problem with the general Tikhonov regularization (TRTLS) is a non-convex optimization problem with local non-global minimizers.Xia [12] proposed an efficient branch-and-bound algorithm (algorithm BTD) for solving TRTLS problems guaranteed to find a globalϵ-approximation solution in most O(1/ϵ) iterations, and the computational effort in each iteration is O n 3 log(1/ϵ) .Beamforming is one of the most important techniques for enhancing the quality of signal in array sensor signal processing, and the performance of a beamformer is usually related to the design of array configuration and beamformer weight.In [13], Chen first proposed a design model for a proximal sparse beamformer, which obtains sparse and robust filter coefficients by solving the composite optimization problem.The objective function of the model is the sum of the least squares term, the approximate term, and the ℓ 1 -regularization term.
Hansen often uses generalized singular value decomposition (GSVD) to analyze regularization methods [14].However, using the GSVD method to solve large-scale discrete ill-conditioned problems requires a large amount of computation and memory requirement.For this kind of problem, Martin and Reichel [9] proposed a method to find the corresponding truncated regularization (TR) solution by using low-rank partial singular value decomposition.In order to improve the time complexity, this paper uses GSVD technology to deal with Tikhonov regularization TLS and establishes Tikhonov regularization TLS based on the GSVD (GTRTLS) algorithm.At the same time, for medium-and large-scale problems, in order to reduce the storage requirements and accelerate the speed of GSVD, the randomized GSVD method [15,16] is used, and then we obtain the randomized GTRTLS (RGTRTLS) algorithm.For the randomized algorithms of large-scale matrix decompositions and their application to ill-conditioned problems, one can see [17][18][19][20] for examples and details.
Our main contribution is to use GSVD technology to deal with Tikhonov regularization TLS (GTRTLS) and to adopt the randomized techniques of [15,16] to implement the GTRTLS procedure in the regularization.The randomized GSVD requires much less storage and computational time than the classical schemes.Numerical examples show the effectiveness and superiority of our algorithms.
This paper is organized as follows: Section 2 describes our technique of combining Tikhonov regularized TLS and GSVD.Section 3 contains our randomized algorithms, and their error analyses for randomized algorithms are in Section 4. The improvement in time and memory requirements is illustrated with numerical examples in Section 5. Section 6 concludes this paper.
Tikhonov Regularization TLS and GSVD
The regularized TLS problem can be expressed as where δ is a positive constant.Typical examples of the matrix L are the first derivative approximation L 1 and the second derivative approximation L 2 , which are as follows (see [14], Equation (1.2), and [21], Equation (4.57)): More precisely, derivative-based finite-difference methods L 1 and L 2 are approximations of the first and second derivative operators on a uniform grid, where the scaling factor is ignored.
The corresponding Lagrange multiplier formulation is where µ is the Lagrange multiplier.
To ensure that the TRTLS problem (1) has a unique solution, throughout this paper, we assume that where K ∈ R n×s is a matrix whose columns form an orthonormal basis of the null-space of the regularization matrix L, and σ min denotes the minimal singular value of its argument [4].
A popular approach to overcoming numerical instability is Tikhonov regularization TLS [4].It can be seen that the regularized total least squares solution can be obtained from the following theorem (see reference [7]): Theorem 1 ([7]).With the inequality constraint replaced by equality, the TRTLS solution x to (1) is a solution to the problem where the parameters λ I and λ L are given by , and where µ is the Lagrange multiplier in (2).λ I and λ L are related by Moreover, the TLS residual satisfies For problem (1), we have the following assumptions: According to the literature [1,12,22], the GSVD of matrix for {A, L} is where U ∈ R m×m and V ∈ R p×p are orthonormal matrices, X is an invertible matrix.The matrices It can be seen that ( 4) is equivalent to the augmented system In order to improve the time complexity, this paper uses GSVD technology to deal with Tikhonov regularization TLS.In the first step, we reduce the GSVD of {A, L} to (7) Let U = U 1 , U 2 ; then, we have or In the second step, using the Elden algorithm [4], only p steps of Givens transformation are needed to eliminate the λ L 1/2 M 0 , which can be expressed as When G is applied to the augmented system (9), we have Since the solution of ŝ can be obtained from the above formula, only the following system can be considered: In the third step, ΣX −1 is reduced to n × n bidiagonal matrix B by orthogonal transformation, such that Finally, through a series of Givens transformations, the above system, whose coefficient matrix can be transformed into a 2n × 2n symmetric indefinite tridiagonal matrix, can be solved by Gaussian partial principal component selection strategy.
To sum up, we call the above algorithm a Tikhonov regularized total least squares algorithm using GSVD technology.It is called Tikhonov regularized total least squares based on generalized singular value decomposition (GTRTLS algorithm for short).
Remark 1.In order to overcome the ill-posedness, we can discard the element close to 0 of item Σ in GSVD, that is, truncated GSVD (TGSVD) and L (see Equation ( 6)), where Σ p (n−k) (n − p ≤ k ≤ n) equals Σ p with the smallest n − k σ i 's being replaced by zeros.In TGSVD, the main information of the original system is retained by choosing the appropriate parameter k, and then the truncated system is obtained by the truncation regularization method.In other words, we combine truncated GSVD and TR to achieve a better regularization effect, which is called TGTRTLS; the expression is as follows: Remark 2. According to Theorem 1, combined with Formula (10), the values of parameters λ I and λ L can be given more effectively.Statistical aspects of a negative regularization parameter in Tikhonov's method are discussed in [7].
Randomized GTRTLS Algorithms
In recent years, there have been many research results of randomized algorithms [15,16].In the truncated case, the randomized algorithm can take the subspace as a random sample and capture most of the information of the matrix, that is, a large-scale problem can be randomly projected into a smaller subspace and also contain its main information, and then some regularization methods are used to solve the small-scale problem.In particular, for severe ill-conditioned problems, we find that GSVD combined with the randomized algorithm is more effective than the classical GSVD method.The general idea is as follows: First of all, with high probability, one can select an orthonormal matrix Q ∈ R m×(k+s) such that ∥A − QQ T A∥ ≤ cσ k+1 , where σ k+1 is the (k + 1)-th largest singular value of A, and c is a constant which depends on k and the oversampling parameter s.It satisfies that R(A T Q) ⊆ R(A T ); here, R(A T Q) is the approximate subspace spanned by the dominant right singular vectors of A. Next, a matrix ((Q T A) T , L T ) T with a small scaled is obtained which can be used to calculate the GSVD of (A T , L T ) T , approximately where U ∈ R m×l and V ∈ R p×p are orthonormal, Z = X −1 ∈ R n×n is nonsingular, and C ∈ R l×n and S ∈ R p×n are rectangular diagonal matrices.Randomized sampling can be used to identify a subspace that captures most of the action of a matrix [15].It provides us with an efficient way for truncation.A large-scale problem is projected randomly to a smaller subspace that contains the main information; then, the resulting small-scale problem can be solved by some regularization methods.Especially for severely ill-posed problems, randomized algorithms are much more efficient than the classical GSVD.So, the advantage of this algorithm is obvious when m ≫ n.The detailed implementation process is shown in reference [16].For the convenience of reading, we describe it as follows (Algorithm 1): Now, we use randomized GSVD technology to deal with Tikhonov regularization TLS.In the first step, the approximate augmented system of augmented system (7) can be obtained by using randomized GSVD and we have In the second step, we use Givens transformation to eliminate λ L 1/2 S, which can be expressed as ; when G is applied to the augmentation system, we can get The solution of V 2 T s can be obtained from the above equation, so only the following system can be considered: In the third step, ΣX −1 is reduced to bidiagonal matrix B by orthogonal transformation such that Finally, through a series of Givens transformations, the above system, whose coefficient matrix can be transformed into a 2n × 2n symmetric indefinite tridiagonal matrix, can be solved by Gaussian partial principal component selection strategy.
To sum up, we call the above algorithm a Tikhonov regularized total least squares algorithm using randomized GSVD technology.It is called Tikhonov regularized total least squares based on randomized generalized singular value decomposition (RGTRTLS algorithm for short).
Error Analysis for Randomized Algorithms
First, we would like to review an important result of [16] regarding randomization algorithms.
Lemma 1 (see [15], Corollary 10.9).Suppose that A ∈ R m×n has the singular values Let G be an n × ( k+s) standard Gaussian matrix with k + s ≤ min{m, n} and s ≥ 4, and let Q be an orthonormal basis for the range of the sampled matrix AG.Then, with a probability that is not less than 1 − 3s −s .
Next, a basic theory of perturbation analysis for TRTLS problems is needed.
Theorem 2 ([10]
).Consider the TRTLS problem (1) and assume that the genericity condition σ n (A)> σ n+1 ((A, b) ) holds.If ∥(δA, δb)∥ F is sufficiently small, then we find that where Next, Lemma 1 is applied to the regularization system (1).Since the randomized system (11) can be seen as its perturbation, the following theorem is obtained from Theorem 2.
Theorem 3. Let
σ n be the singular values of matrix A, and α = c/∥A∥ 2 , c= 1 + 6 (k + s)logs + 3 (k + s)(n − k), with the matrix A T , L T T as in (6).Suppose that Algorithm 1 is executed using the Gaussian matrix G to achieve GSVD approximation of matrix pairs A T , L T T .Assumption (5) is satisfied.x trtls is the solution of (1), and x gtrtls is the minimum two-norm solution of the problem (11), δx = x trtls -x gtrtls .Then we have with a probability greater than 1 − 3s −s .
Numerical Examples
In this part, we illustrate the effectiveness and superiority of our methods through specific examples.We use the regularization tool package to perform the calculation on MATLAB R2016a.
Example 1.
The test problem is obtained by executing the function ilaplace (n, 2).The matrix A and the exact solution x are given such that ∥A∥ F = ∥Ax∥ 2 = 1, and the perturbed right-hand side is generated as b = (A + σ∥E∥ −1 F E)x + σ∥e∥ −1 2 e, where the perturbations E and e are formed by a normal distribution with zero mean and unit standard deviation.L is the first derivative operator.The dimensions are m = n = 39.Noise levels are taken as σ = 0.001, σ = 0.01, σ = 0.1 σ = 1.
We see that for small values of σ and for the same value of λ L , the three methods result in almost identical minimum relative errors.However, for a larger value of σ, the minimum relative errors of the GTRTLS method and the RGTRTLS method are significantly smaller than that of the TRTLS method, and they occur at smaller values of λ L as shown in Table 1 and Figure 1.So, the potential advantages of the GTRTLS method and the RGTRTLS method are shown.with a probability greater than 1 − 3 .
Numerical Examples
In this part, we illustrate the effectiveness and superiority of our methods through specific examples.We use the regularization tool package to perform the calculation on MATLAB R2016a.We see that for small values of and for the same value of λ , the three methods result in almost identical minimum relative errors.However, for a larger value of , the minimum relative errors of the GTRTLS method and the RGTRTLS method are significantly smaller than that of the TRTLS method, and they occur at smaller values of λ as shown in Table 1 and Figure 1.So, the potential advantages of the GTRTLS method and the RGTRTLS method are shown.We find that the calculation time of the RGTRTLS method is less than that of the GTRTLS method, and the GTRTLS method is less than that of the TRTLS method, as shown in Table 2. where τ = 2 r and (m, 1)−1.It is easy to verify that ∥b − b∥/ ∥b∥ = σ.We set σ = 0.001 and the size n = 1024 in the tests.The matrix L is L 1 and the regularization parameter λ L and λ I are selected by Remark 2.
For a better understanding of the tables below, we list the notation here: • x is the true solution of the TLS problem (1).
•
x gsvd is the solution of (1) by classical GSVD.
•
x rgsvd is the approximate regularized TLS solution in (11) by randomized algorithms.
For n = 1024, the corresponding errors and time are shown in Table 3, and the performance is shown in Figure 2. We apply the GTRTLS algorithm and the RGTRTLS algorithm to Example 2 and compare the errors and execution times.The randomized approach in Algorithm 1 still shows good performance in Table 3 and is competitive compared with the classical GSVD, judging from the errors E gsvd and E rgsvd and the execution times t gsvd and t 1gsvd .
We cannot solve large-scale or more complex ill-conditioned problems, such as n = 4096, using classical SVD or GSVD due to the high memory requirements.So, one can use preconditioned techniques first, and then use our method for computation.We apply the GTRTLS algorithm and the RGTRTLS algorithm to Example 2 and compare the errors and execution times.The randomized approach in Algorithm 1 still shows good performance in Table 3 and is competitive compared with the classical GSVD, judging from the errors and and the execution times and .We cannot solve large-scale or more complex ill-conditioned problems, such as n = 4096, using classical SVD or GSVD due to the high memory requirements.So, one can use preconditioned techniques first, and then use our method for computation.
Conclusions
In this paper, the generalized singular value decomposition technique is used to deal with Tikhonov regularized total least squares problems to approximate the true regularized TLS solutions, and the GTRTLS algorithm is proposed.The time complexity of the GTRTLS algorithm is better than TRTLS proposed by Golub, Hansen, and O'Leary.For medium-and large-scale problems, in order to reduce the storage requirements and accelerate the speed of GSVD, this paper adopts the random GSVD method and obtains the RGTRTLS algorithm.Numerical examples show that our algorithm has obvious effectiveness and superiority.
2 : 5 : 6 :
where n − p < l < min{m, n}.Output: Orthonormal U ∈ R m×l and V ∈ R p×p , rectangular diagonal C ∈ R l×n and S ∈ R p×n and nonsingular Z = X −1 ∈ R n×n .1: Generate an n × l Gaussian random matrix Ω; Form the m × l matrix Y = AΩ; 3: Compute the m × l orthonormal matrix Q via the QR factorization Y = QR; 4: Form the l × n matrix B = Q T A; Compute the GSVD of {B, L} in (Form the matrix U ∈ R m×l , U = QW and denote Z = X −1 in(12).
Example 1 .
The test problem is obtained by executing the function ilaplace (n, 2).The matrix A and the exact solution x are given such that ‖‖ = ‖‖ = 1, and the perturbed right-hand side is generated as = ( + ‖‖ ) + ‖‖ , where the perturbations E and e are formed by a normal distribution with zero mean and unit standard deviation.L is the first derivative operator.The dimensions are = = 39.Noise levels are taken as = 0.001, = 0.01, = 0.1 = 1.
Figure 1 .
Figure 1.Exact solutions, TRTLS solutions, GTRTLS solutions, and RGTRTLS solutions under four values of the noise levels σ.
Figure 2 .
Figure 2. The comparison for exact solutions, GTRTLS solutions, and RGTRTLS solutions under the value of the noise levels = 0.001.
Figure 2 .
Figure 2. The comparison for exact solutions, GTRTLS solutions, and RGTRTLS solutions under the value of the noise levels σ = 0.001.
Table 2 .
Time comparison table of TRTLS method, GTRTLS method, and RGTRTLS method.
Table 3 .
The comparison of GTRTLS method and RGTRTLS method.
Table 3 .
The comparison of GTRTLS method and RGTRTLS method. | 4,405.8 | 2024-06-21T00:00:00.000 | [
"Computer Science",
"Mathematics"
] |
Research on Influence of Manager ’ s Innovation Preference on Innovation-Decision Making
It is one of the key factors which cause “innovation dilemma” that managers prefer to support the sustaining innovation project. From the view of the manager’s innovation preference, the main propose of the paper is to study why it happened. The manager’s innovation preference will guide and motivate the staffs how to innovate, therefore it is appropriate to analyze it by using the principal agent theory. Conclusions can be got by establishing and analyzing a multi-task principal-agent model. First of all, the model basically explains why incumbent enterprises prefer adopting sustaining innovation and entrant enterprises are inclined to disruptive innovation project. Secondly, the selection rights of middle managers towards innovation projects determine the strategic direction of enterprises. Manager’s innovation preference is consistent with the innovation types of employees. At last, the paper suggests that incumbent enterprises should indeed establish self-organizations or spin-off organizations to better carry out disruptive business.
Introduction
When a manager faces multiple ideas, he definitely needs choose some of them and then turn them into formal innovation projects.Because the resources are limited, in practice, middle managers are usually responsible for doing these works, and then the result will be delivered to senior managers after selection.Generally, middle managers will subconsciously avoid proposing an idea which may not be approved by top manager.Obviously the top manager's innovation preference is on behalf of an enterprise's innovative strategy.So if the middle manager's innovation preference is consistent with top manager That is to say, when middle managers can determine which project can become formal innovative project with strategic value to enterprise, to some extent, the middle managers guide the strategic direction of the enterprise.
Over the years, the middle managers haven't got enough attention on enterprise innovation.Nanaka (1995) sharply points out that middle managers should never be "disappearing level", they are the intersection of vertical and horizontal information flows and of great significance in the process of organizing knowledge creation [1].Christensen (2003) also agrees that middle managers play a key role in creativity stereo-type process [2].From the perspective of disruptive innovation, the innovation strategy of enterprise which focuses on the core business that enterprise input all resources support, is called sustaining innovation strategy, instead of targeting at consumers of the main market, disruptive innovation strategy is more inclined to new consumption market or even low consumption market.Why incumbent enterprises prefer sustaining innovation compared with disruptive innovation favored by new enterprises?Some scholars argued that the reason is incumbent enterprises have formed a relatively stable core conventional business, which has a competitive advantage, especially core competitive advantage [3], their managers need to pay more attention to maintain it [4].The sustaining innovation strategy aiming at core customers and because sustaining innovation won't change the existing value net-works of the enterprise, therefore its risks are relatively controllable and profits predictable.Thus, incumbent enterprises prefer supporting sustaining innovation projects.Comparatively, entrant enterprises haven't formed their core business.The limited profits from core business push them to encourage innovation.Under the circumstances of opening creative minds and tolerant for failure, the idea proposed by employees to face new consumption market or even low consumption market will easily be adopted.That is to say, the distinguished features of diverse enterprises determine that the middle managers in incumbent enterprises prefer sustaining innovation projects, and the managers of entrant enterprises like disrup-tive innovation projects.
Core advantage often means "core rigidity" [5].Under the guidance of sustaining innovation strategy, an enterprise's innovations focus on high-end market, thus innovation consequence may be highly concentrated on specific areas.The innovation scope seems too narrow in the technological era, which is a sign of lacking innovations in essence.Christensen (2003) thought that the overemphasis on core customers would probably result in the ignorance of disruptive opportunities.When the products performance overshoots over time, core customers would be reluctant to pay for the performance premium.On the other hand, in the early stage, although the produce of entrant enterprises could not meet the requests of core customers, they keep improving products performance and ultimately deprive market of incumbent enterprise.This is the key reason why some enterprises run well and allocate resources with profits maximization theory but still lost market competitiveness.
In reality, the managers in incumbent enterprises have negative attitude towards even oppose the disruptive innovation projects innately, which results in strategic failures, such as Kodak bankruptcy.Kodak devoted itself to various imaging technique researches and was equipped with strong ability to develop and promote various imaging techniques, but the vast resources put in core film business by managers results in the unfavorable position of digital imaging techniques research and promotion.As a consequence, when a series of digital products appear in the market, Kodak declined rapidly, ending up with bankruptcy.
The structure of this paper is: After introduction, from the view of the innovation preference of middle managers, a principal-agent model will be established, then comparative static analysis would be done.After that, the author thinks that further discussions are necessary based on human capital theory and dynamic capability.The last is the section of conclusions.
Hypothesis and Solution of the Model
The paper names sustaining innovation activity and disruptive innovation activity as conventional business and innovative business.It assumes that the choices of middle managers towards innovative projects are consistent with those of supreme decision-making level.Managers are principals and employees are agents.Under the standard principal-agent model, agents have two specific tasks: conventional task and innovative task.The following hypothesis can be got based on the above analysis.Firstly, the standard multi-task principal-agent model assumes that principal and agent sign a contract based on performance: S = α + β 1 x 1 , α represents the fixed salary of agent; β 1 represents the sharing coefficient produced by the accomplishment of the conventional task or incen-tive factor given by principal to agent, x 1 is the output that can be confirmed.Then x 1 = e 1 + ε 1 , e 1 shows the efforts agent spends in finishing conventional task, ε 1 is random variable which is independent from the efforts of the agent ε 1 → (0, σ 2 ).Both principal and agent are riskneutral, the reservation utility of agent is zero.Further hypotheses on this basis: Hypothesis 1 Assume that the principal innovation preference is continuous and represented by parameter k (0 < k ≤ 1).The principal is reluctant to invest resources in developing innovation products when k is large, and willing to put existing resources into core business.The principal is cautious when get involved in new business, thus those ideas concerned with traditional core business are more easily get support from principal, using existing resources and knowledge to finish conventional core business.The sense of innovation of the principal is stronger when k is smaller.The principal supports noncore business ideas, having strong desire to explore new market, motivating agent to use new knowledge and paying more attention to the expected profits from the accomplished innovative tasks.
Hypothesis 2 Principal and agent are players in a game, the time sequence of the game is: the first stage, principal has innovation preference, which is reflected on the expected profits of the principal.Then there is an explicit incentive contract, focusing on conventional task, between principal and agent.The second stage, agent has an idea of facing new market and exploring new customers, principal considers whether to change the idea into a formal innovative project according to his own orientation.The support of the principal means that agent has to decentralize efforts to improve the innovation idea under the premise of having original schedule tasks.Thus, the choice of the principal would lead the energy allocation of the agent.The third stage, according to the performance signal, principal pays salary to the agent according to the contract.If there are innovation achievements, the two parties proportionally allocate them according to the innovation value.
Hypothesis 3 Changing innovative idea into an innovation achievement needs agent's efforts, innovative efforts, represented by e 2 , unit effort cost is c 2 , and e 1 is conventional effort, the corresponding unit effort cost is c 1 , and e 1 + e 2 = 1.If principal wants to put his idea to practice, he must face the choice between innovative effort and conventional effort.When agent implements e 2 and gets innovative achievement, the achievement value is: μe 2 , m represents the innovative ability parameter (0 ≤ μ ≤ 1), or marginal output of the innovative efforts.The results can be observed, but cannot be confirmed by the third party.The allocation ratio is λ (0 ≤ λ ≤ 1).Then the expected innovative profits of principal and agent are λμe 2 and (1 -λ)μe 2 respectively.
Hypothesis 1 is the central hypothesis of this article.Hypotheses 2 and 3 are about innovation achievement sharing, Aghion and Tirole (1994) assume that there can be renegotiation after innovation achievement [6].The result of the renegotiation can be a Nash Equilibrium Solution.In the case of repeated game, agents can perform a ruthless strategy to ensure negotiations.The negotiating capacity of the agent is determined by whether the essential human resource capital is general or specific.Base on this, from the view of the appropriable ability of innovation, (Helleman 2011) studies the incentive problem [7].
Under the above hypotheses, the profits of agent's conventional efforts and innovative efforts are α + β 1 e 1 and (1 -λ)μe 2 the efforts are 2 2 c e , c 2 e 2 respectively.The agent's certain earning in the third stage is The derivation of U A (e 1 , e 2 ) to e 1 : Under the previous circumstances, take principal's innovation preference into consideration and combine agent's incentive constraints and participation constraints, the expected profits of the principal in the second stage is: The principal's decision in the first stage is: We can assume that V stands for total certain earnings, put IR into the objective function: 1 k e Then put IC into the objective function, and the derivation is: After simple calculation, from the principal's objective function, we can get that the second derivative of β 1 is 1 , and the second order derivative of e 1 , e 2 is -c 1 < 0. Thus, the objective function V of principal is the strictly concave function of β 1 , e 1 , e 2 .
The Model Analysis nce of Principal's
Prop vious hypotheses, in the first
Analysis on the Influe Innovation Preference on Innovation Decision Making
osition 1 Under the pre stage of the game, principal can choose to support traditional innovation, principal pays attention to the profits brought by conventional task (k → 1).He can also choose those innovations facing new market and customers, that is to say, principal focuses on the profits brought by innovative task (k → 0).Under the equilibrium, the optimal value of objective function was got at the peak (k → 0 or k = 1).Besides, the target expected profits of innovative task efforts are higher than those of conventional efforts So V * is a strictly concave function of k.That is to say, the optimal value of the objective function can be got at the peak point.Where exists a k * , when .ns that prefere This mea nce coefficient (k) will enter th : e upward moving channel (k → 1), the principal tends to support the tasks related with traditional core business.When k * > k, the preference coefficient (k) will enter the downward moving channel, the principal tends to focus on innovations facing new market and customers; he cares about the profits brought by new innovation tasks.From (5) we can know that this task is risky and at the expense of losing traditional business profits.
At last, when k = 1, the objective function is Similarly, when k → 0, the objective function is: at Proposition 1 points out th the principal is not quite willing to innovate new business under certain circum-stances (k * < k), he likes to use the existing knowledge and experience to engage in less risky innovation.When (k * > k), the innovation is risky because the innovation result is uncertain.
Under the condition that innovation preferences of pri es new ma ther parameters unch ing fu ncipal and enterprise are consistent, when k = 1, the principal type is named the optimal enterprise which allocates resources according to profits maximization [8].When conventional core business relatively stabilizes, enterprise naturally orients its resources on conventional core business, which is reflected by k = 1.As to innovative strategies, the optimal enterprise also innovates and chooses conservative strategies.When k → 0, the principal type is called innovative enterprise, which focuses on the development opportunities in the future and stands for the profits loss in short period, motivating innovation and tolerating failure.In reality, the entrant enterprise hasn't founded any core business, the profits from which are limited.Thus, encouraging staff innovation can bring more profits and aggressive innovative strategy is more likely to be adopted, such as disruptive innovation.k → 0 reflects this feature of newly established enterprise.The value of the innovation preference (k) proves that traditional incumbent enterprises like to focus on conventional core business, while newly established enterprise is willing to allot no effort in innovation business.From , we can know that the innovation which facrket will be more risky, but it can bring more profits once it succeeds.This can briefly explain one question in disruptive innovation theory: why incumbent enterprise likes sustaining innovation, and why entrant enterprise prefers disruptive innovation.The lower the innovation preference is, the more the principal is willing to maintain core customers, which can be regarded as sustaining innovative strategy.The higher the innovation preference is, the more the principal is willing to adopt the innovative strategy facing new market and customers, which can be considered as disruptive strategy.As a principal, the innovative preference of middle managers influences the innovation decision-making and finally affects the strategic direction of the enterprise.
Proposition 2 Suppose all the o anged under the previous assumption, when the preference coefficient (k) enters the upward channel (k → 1) the motivation of conventional business will be strengthened.The principal's innovative efforts will decrease continuously if he gets less encouragement.On the contrary, the energy spent on conventional business will increase.When the principal totally tolerates innovation (k → 0), the optimal dominant performance motivation coefficient is negative, and the agent will assign no effort in innovative business.When the principal cannot tolerate the innovation of the agent (k = 1), the optimal dominant performance coefficient reaches the maximum, thus attract-ing agent to devote more in traditional core business.
Proof: From ( 6), we can know that k is an increas nction of 1 And when k = 1, 1 1 . From (7), from the assumption e 1 ≥ 0, so The proposition 2 points out that when the principal's innovative intention is weak (k → 0), its incentive towards traditional business is bigger compared with strong (k → 0).The combination of proposition 1 and 2 shows when the principal completely supports the agent's innovation activity, the agent will get motivated and devotes all his energy to innovative task.When the principal doesn't support agent's innovation at all (k = 1), the efforts of agent on traditional business is , and is left to be used on innovation.The re, the on preference greatly influences the devotion degree of the agent, and further influences the output efficiency of the innovation.
Lemma 1 Assuming that the refo innovati re are two types of agents: th in (3), from the previous as e conventional type lacking creative spirits, or the aggressive innovative type.When the principal doesn't know which type the agent belongs to, he can provide contract sets {(α, 0), (0, β)} for agent to choose, which reflects the type of the agent.
Proof: Put ( 6) and ( 7) into IR sumption , we can get the optimal fixed income: From Lemma 1, we can know that α * decrease s k incr s a eases.While from Proposition 1, the objective function gets the optimal value at the end point (k → 0 or k = 1).Thus, when the orientation of the principal is conservative (k = 1), principal will pay more attention to traditional business profits.At this moment, continuously reducing fixed wage is the optimal choice for principal.Because of the constraints of limited liability, α ≥ 0, as k increases, α * = 0 at last.When the principal wants to innovate (k → 0), from Proposition 2, the optimal 1 is Copyright © 2012 SciRes.TI negative.When k → 0. At this moment, the conventional task is not motivated; the profits of the principal and agent depend on the agent's innovative efforts (e 2 ).Because of the uncertainty of the innovation, the principal will not determine the profits of the agent solely by innovative achievements sharing coefficient (λ), thus there must be α * > 0. If assume that β ≥ 0 reasonably, then the optimal explicit contract given by principal to agent is (α, 0).When the principal's orientation is conservative (k = 1), the optimal explicit contract given by the principal to the agent is (0, β).In reality, the enterprise recruitment pays more attention to whether the employees' values are in consistent with the enterprises'.Employees who are willing to take risks are welcomed by enterprises encouraging innovation.Through different compensation designings, enterprises can easily find employees meeting its own culture.As the enterprise grows, it can get more and more innovative employees together.It is the opposite case for enterprises under the traditional culture.
Prop
Under the previous assumptions, suppose t
Sharing Coefficient on Effort Level and Profits
osition 3 all the other parameters unchanged, as λ increases, the agent is more reluctant to engage in innovative activities, the effort level of the innovative task will continuously decrease, the motivation of conventional business is strengthened.On the contrary, the effort level of conventional task will increase constantly.
Proof: from (6), it is obvious tha 1 is the increasing function of λ.From (7) we can know that According to the assumption, willing to motivate the agent to allot more innovative efforts, the less the agent gets profits from the innovative achievement.Thus, agent likes to engage in conventional activity instead of focusing on innovative efforts.Especially, when λ = 0 (or λ = 1), the agent (or the principal) gets all the innovative achievements.When λ = 0, agent gain all the returns, principal has nothing; when λ = 1, the principal got everything, and agent has nothing and finally gives up innovative efforts.Thus the optimal λ must be between (0, 1).
Proposition 1 to 3 say that principal can interfere the more the principal i with th suppose all the other parameters unch s e innovation of agent; it can also lead the energy allocation by influencing the explicit motivation strength.Of course, on the basis of this, the energy allocation of agent is also determined by the of the principal's appropriablity from innovation.
Proposition 4 anged, there exists a threshold value λ under equilibrium condition, when λ * < λ, the objective function constantly increases as λ increases.When λ < λ * , the objective function will constantly decreases as λ increases.
Proof: from (8) s, the objective V * is the u-curve of λ.Therefore, th Thu ere must exist a λ * , when λ < λ * , 0 Pr sh oposition 4 ows the influence of innovation result ex
Further Discussions
tal theory, the principal's clusivity (λ) on (V * ).For certain innovative efforts, it is obvious that the stronger principal's appropriable ability, the more profits he can get; the principal thus have strong motivation to encourage agent to innovate.But from proposition 3, when λ → 1, agent will not innovate at all, thus the innovation profits declines to 0. Therefore, when principal's appropriable ability is strong, agent would not like to engage in innovation.On the contrary, when λ → 0, from proposition 3, if agent likes to innovate, but principal cannot get any innovative profits, thus, the weaker principal's appropriable ability, the more reluctant for the principal to motivate the agent to innovate.So on the topic of innovation results sharing, proposition 4 gives the direct proof of the conflicts in innovation fruits allocation, both parties need to compromise to reach a balance.
From the view of human capi negotiation ability depends on whether the human capital is general or specific.Specific human capital plays a key role in transaction cost theory [9] and property rights theory [10].On the other hand, if agent's innovation activity requires the principal's specific assets, then this is the so-called appropriability problem.Nelson (1959), Arrow (1962), Teece (1986) had analyzed the problem from different views [11,12].The PFI theory proposed by Teece provided a new view for analyzing appropriability problem [13].One important concept in PFI theory is Complementary Assets, which is a key factor influencing enterprise's innovative profits exclusivity.The strong complementary assets of an enterprise mean its strong ability to get profits from innovation.In the structure of this paper, the bargaining power of both principal and agent on innovation profits allocation depends on the human capital ability of the agent and the principal's appropriable ability from innovation.From the view of complementary assets, combine the innovation preference coefficient, proposition 4 says when the principal's innovative preference is high (k → 0); he is concerned with profits coming from innovative business.At this moment, if agent's innovation must depend on the complementary assets of the principal, then the incentive on We assume the preference is in accordance with innovation strategy of enterprise, from the view of the manager's innovation preference, the author established a multi-task principal-agent model that explains how the manager's innovation-decision to incentive the staffs for what kind of innovation.After analyzing the theoretic model, conclusions can be got as follows: Under the framework of model, we can g in conditions, the selection rights of middle managers towards innovation projects determine the strategic direction of enterprises.The model briefly explains why incumbent enterprises like so-called sustaining innovation and entrant enterprises are inclined to disruptive innovation.
According to manager's innovation preference, enterse can design different employment contracts.The combination of different contracts can not only reflect manager's innovation preference but also show employees' innovation types.Providing different contract sets for employees, the type of risk avoidance employees would like to choose enterprise with innovative culture, and the type of risk avoidance tends to choose traditional enterprise.
Under the condition the principal's appropriable abilto develop disruptive innovation project is not strong, if the manager's innovation preference is inclined to sustaining innovation, enterprise won't encourage employees to innovate disruptive projects, in this way employees either choose to leave, or give up disruptive projects and only focus on conventional task.When the manager's innovation preference is more like sustaining innovation, if enterprise wants to share employees' innovation achievement, he needs to play the role just as venture capitalist.This shows that incumbent enterprise should indeed establish self-organization or spin off organization to better finish disruptive task.
Regarding the innovation projec esn't influence the conclusions, but this paper analyzes the issue from individual aspect, neglecting the efforts conflicts among agents, which needs further research.
The research is financed by Foundation of China.No: 71172095, MOST of China Science and Technology Basic Tasks FANEDD No: 2011IM020100.
's explicit performance is so-called low-level motivation.The principal can ask for high proportion of innovation profits allocation.The principal needs to reduce the sharing coefficient of the innovative returns to motivate the agent.On the contrary, because it is incredible to give agent high proportion of innovation results sharing, so the agent would probably decide to resign if agent's innovation doesn't need the principal's complementary assets.The abundant evidence shows that the new enterprises engaging in disruptive business are constituted by turnover workers from incumbent enterprises.Thus, from the view of disruptive innovative theory, if principal still cares for the innovation achievement from agent, he needs to consider whether establish self-organization or spin off organization (Christencen 2003), making agent engage in disruptive business flexibly and eager to get innovative fruits.Otherwise, agent would leave. agentth | 5,597.2 | 2012-08-31T00:00:00.000 | [
"Business",
"Economics"
] |
Interplay of ortho- with spiro-cyclisation during iminyl radical closures onto arenes and heteroarenes
Summary Sensitised photolyses of ethoxycarbonyl oximes of aromatic and heteroaromatic ketones yielded iminyl radicals, which were characterised by EPR spectroscopy. Iminyls with suitably placed arene or heteroarene acceptors underwent cyclisations yielding phenanthridine-type products from ortho-additions. For benzofuran and benzothiophene acceptors, spiro-cyclisation predominated at low temperatures, but thermodynamic control ensured ortho-products, benzofuro- or benzothieno-isoquinolines, formed at higher temperatures. Estimates by steady-state kinetic EPR established that iminyl radical cyclisations onto aromatics took place about an order of magnitude more slowly than prototypical C-centred radicals. The cyclisation energetics were investigated by DFT computations, which gave insights into factors influencing the two cyclisation modes.
General experimental section
All reagents and solvents were purchased from either Sigma Aldrich or Alfa Aesar and used without further purification. Toluene and tetrahydrofuran were distilled over sodium, and dichloromethane was distilled over calcium hydride. Benzaldehyde oxime and acetophenone oxime were prepared according to the literature procedure [1], as was N-benzylpent-4-en-1amine [2]. Column chromatography was carried out using Silica 60A (particle size 40-63 µm, Silicycle, Canada) as the stationary phase, and TLC was performed on precoated silica gel plates (0.20 mm thick, Sil G UV 254 , Macherey-Nagel, Germany) and observed under UV light. 1 H and 13 C NMR spectra were recorded on Bruker AV III 500, Bruker AV II 400 and Bruker AV 300 instruments. Chemical shifts are reported in parts per million (ppm) from low to high frequency and referenced to the residual solvent resonance. Coupling constants (J) are reported in hertz (Hz). Standard abbreviations indicating multiplicity were used as follows: s = singlet, d = doublet, t = triplet, dd = double doublet, q = quartet, m = multiplet, b = broad. Melting points (mp) were determined using a Sanyo Gallenkamp apparatus and are reported uncorrected. Mass spectrometry was carried out at the EPSRC National Mass Spectrometry Service Centre, Swansea, UK.
Synthesis and experimental section
Oxime carbonates 1a-f, and 2a,b were prepared as described previously [1].
UV cyclisation of oxime carbonate derivatives general procedure
A quartz tube was charged with oxime carbonate (1.0 equiv), 4-methoxyacetophenone (MAP) (1 equiv wt/wt) and benzotrifluoride (3 mL). The reaction mixture was degassed by bubbling Ar through the solution for 15 min. The solution was irradiated with UV light (400 W medium pressure Hg lamp) for 3 h. The solvent was removed under reduced pressure and the crude residue purified by column chromatography (CH 2 Cl 2 /EtOAc 9:1 as eluent).
EPR spectroscopy
EPR spectra were obtained at 9.5 GHz with 100 kHz modulation employing a Bruker EMX 10/12 spectrometer fitted with a rectangular ER4122 SP resonant cavity and a Bruker ER4122-SHQE X band cavity on EMX and EMX Micro consoles in Manchester. Stock solutions of each oxime carbonate (2 to 15 mg) and MAP (1 equiv wt/wt) in tertbutylbenzene or benzene (0.5 mL) were prepared and sonicated where necessary. An aliquot (0.2 mL), to which any additional reactant had been added, was placed in a 4 mm o.d. quartz tube and deaerated by bubbling nitrogen for 15 min. Photolysis in the resonant cavity was by S8 unfiltered light from a 500 W super pressure mercury arc lamp or, in the Manchester experiments, the light source was a Luxtel CL300BUV lamp. Solutions in cyclopropane were prepared on a vacuum line by distilling in the cyclopropane, degassing with three freeze-pump-thaw cycles and finally flame sealing the tubes. In all cases where spectra were obtained, hfs were assigned with the aid of computer simulations using the Bruker SimFonia and NIEHS Winsim2002 software packages. For kinetic measurements, precursor samples were used mainly in "single shot" experiments, i.e., new samples were prepared for each temperature and each concentration to minimise sample-depletion effects. EPR signals were digitally filtered and double integrated by using the Bruker WinEPR software and radical concentrations were calculated by reference to the double integral of the signal from a known concentration of the stable radical DPPH [1 10 −3 M in PhMe], run under identical conditions, as described previously. The majority of EPR spectra were recorded with 2.0 mW power, 0.8 G pp modulation intensity, and a gain of ca. 10 6 .
Computational methods
Radical ground-state calculations were carried out by using the Gaussian 09 program package [4]. Becke's three-parameter hybrid exchange potential (B3) was used with the LYP correlation functional, B3LYP. This method has previously described the chemistry of iminyl radicals accurately. The standard split-valence 6-31+G(d) basis set was initially employed and then the computations were extended to the UB3LYP/6-311+D(2d,p) level. Geometries were fully optimised for all model compounds. Optimised structures were characterised as minima or saddle points by frequency calculations. The experimental kinetic and spectroscopic data was all obtained in the nonpolar hydrocarbon solvents tert-butylbenzene or cyclopropane. Solvent effects, particularly differences in solvation between the neutral reactants and neutral transition states, are therefore expected to be minimal. In view of this, no attempt was made to computationally model the effect of the solvent. | 1,159.2 | 2013-06-04T00:00:00.000 | [
"Chemistry"
] |
Exosomal lncRNA SHNG7 Promotes High-Grade Serous Ovarian Cancer Progression and Function As ceRNA to Target Notch1 By Sponging miR-34a-5p
Serous ovarian cancer, especially high-grade serous ovarian cancer (HGSOC) has a high mortality rate, and its five-year survival rate is only 30%. The reason for the difficulty in diagnosis and treatment is that the origin and pathogenic mechanism of HGSOC are still poorly understood. In this study, we tried to explore the molecular mechanism of lncRNA SHNG7 in regulating proliferation, invasion and migration in HGSOC. The expression of lncRNA SHNG7 was upregulated in cancer tissues, and was closely correlated with poor outcomes. LncRNA SHNG7 promoted cell proliferation, invasion, migration and influenced the cell cycle of cancer cell lines. Furthermore, lncRNA SHNG7 could function as competitive endogenous RNA (ceRNA) via directly sponging of microRNA-34a-5p, which further regulating the expression of Notch 1. Moreover, lncRNA SHNG7 could be carried by exosomes and the exosomal lncRNA SHNG7 promoted angiopoiesis. Taken together, our results proved that lncRNA SHNG7 could function as ceRNA and contribute to HGSOC progression, which provided a novel prognostic and therapeutic marker for HGSOC.
Introduction
Ovarian cancer is a common malignant tumor in women, with the fth highest mortality rate. Serous ovarian cancer, especially high-grade serous ovarian cancer (HGSOC) has a high mortality rate (90%), and its ve-year survival rate is only 30% (1). One of the main reasons for the high mortality rate of ovarian cancer is that most patients are in the advanced stage when they are discovered, and metastasis is existed (2). The reason for the di culty in diagnosis and treatment is that the origin and pathogenic mechanism of HGSOC are still poorly understood. In this study, we tried to explore the molecular mechanism in HGSOC.
Long non-coding RNAs (lncRNAs) are a class of transcripts with more than 200 nucleotides in length and no protein-coding capacity. Emerging evidences proved that lncRNAs act as key regulators of target gene expression in various biological processes, such as gene transcription, RNA splicing, and RNA transport and translation (3). And aberrant expression of lncRNAs has been found involved in tumor occurrence, progression and metastasis (4). It is proved that some lncRNAs could function as competing endogenous RNAs (ceRNAs) via sponging microRNAs (miRNAs)to regulate the expression of speci c genes (5). In HGSOC, several lncRNAs have been identi ed as key roles in the initiation or progression. For example, upregulation of lncRNA SOCAR promoted proliferation, migration and invasion in ovarian cancer cells (6). LncRNA NEAT1 was proved to promote cell proliferation and migration by sponging miR-506 in HGSOC (7). While in this study, the function of lncRNA SHNG7 was studied in HGSOC.
Exosomes are a type of membranous vesicles with a diameter of about 30 ~ 100nm, which are released from cells into the extracellular matrix. Many kinds of cells could release exosomes (8). Exosomes could function as a carrier for multiple messenger RNAs, microRNA, lncRNA and circRNA. Exosomes could be isolated from serum, urine, bile and breast milk (9,10). More and more studies have shown that the exosomes of tumor cells are related to the occurrence and deterioration of tumors (11). They can regulate immune function, promote tumor angiogenesis, invasion and metastasis, and even directly affect other tumors or non-tumor cells, thereby affecting the fate of the cell or tissue (12).
In the present study, we found lncRNA SHNG7 was dyregulated in HGSOC tissues, and its overexpression dramatically promoted cell proliferation, migration and invasion. LncRNA SHNG7 could be released from cells via exosomes, and exosomal lncRNA SHNG7 was demonstrated promoting angiogenesis. Besides, we identi ed that lncRNA SHNG7 could function as ceRNA to target Notch1 by competitively sponging miR-34a-5p. These ndings provide new insights into the molecular functions of lncRNA SHNG7 and shed new light on the treatment of HGSOC. and maintained using Dulbecco's modi ed Eagle's medium (Gibco, Carlsbad, CA, USA) supplemented with 10% fetal bovine serum (FBS). All cells were incubated in a humidi ed atmosphere containing 5% CO 2 at 37 ℃. Vectors that stably overexpressed lncRNA SHNG7, siRNA targeting lncRNA and SHNG7, miR-34a-5p mimics, miR-34a-5p inhibitor and controls were transfected into cells using Lipofectamine 2000 reagent (Invitrogen, Carlsbad, CA) according to the manufacturer's instructions.
Quantitative RT-PCR analysis
Total RNA was extracted from tissues and cells using Trizol reagent (Invitrogen) according to the manufacturer's protocol. Extracted RNA was reversely transcribed into cDNA using a PrimeScriptTM RT Reagent Kit (TaKaRa Bio, China) following the manufacturer's instructions. qRT-PCR was carried out using a SYBR Premix Ex Taq™ Kit (Takara). U6 or GAPDH was used as an internal control. The relative expression was determined using the 2 −∆∆CT method.
Cell proliferation
Cell proliferation was detected using Cell Counting Kit-8 (CCK-8; Dojindo, Kumamoto, Japan) according to the manufacturer's protocol. Cells were seeded into 96-well plates. After incubation for the indicated time, measured on a Microplate Reader (Bio-Rad).
Cell cycle assay
Cell cycle analysis was performed using ow cytometry. Cells were collected and wash twice with PBS (1×), and xed with 70% ethanol at -20°C for 24 h. Cells were incubated with RNase A at 37°C for 30 min, and then stained with 400 µL propidium iodide on ice for 30 min. Cell cycle distribution was analyzed using a BD FACSCalibur™ ow cytometer (BD Biosciences).
Transwell assay
Cell migration and invasion were completed with a Transwell system (Corning, NY, USA). A total of 10 4 cells in 200 µL of serumfree medium were added to the upper chamber. For invasion assay, Matrigel was additional coated on the upper chamber. The lower chamber was lled with 600 µl medium with 20% FBS. After incubation, cells on the upper chamber were removed and the cells on the lower surface were xed with 20% methanol and stained with 0.1% crystal violet. The cells were observed under a microscope.
Western blot
Cells were collected and lysed in RIPA buffer containing protease inhibitors. Proteins were separated by 8% SDS-PAGE then transferred onto PVDF membranes (Millipore, USA). Then the membranes were blocked with 5% nonfat milk for 1 h at room temperature, and incubated with primary antibodies overnight at 4°C. Subsequently, the membranes were incubated with secondary antibodies for 1.5 h. the protein bands were detected using the Pro-lighting horseradish peroxidase (HRP) agent. The expression of β-actin was used as the internal control.
Luciferase Assay
The dual-luciferase miRNA target expression vector pmirGLO (Promega, Madison, Wisconsin) was used to generate luciferase reporter constructs. Wild-type (SHNG7-WT and Notch1-WT) and mutant-type (SHNG7-MT and Notch1-MT) vectors were constructed. Cells were cotransfected with miR-34a-5p mimics (or NC) and wild-type (or mutant-type) vectors using Lipofectamine 2000. Luciferase activity was examined using a luciferase reporter assay kit (Transgen Biotech, Beijing, China). Each group was performed in triplicate.
Isolation of exosomes
Exosomes were isolated from cell culture medium by differential centrifugation. Cells and other debris were removed by centrifugation at 300 g and 3,000 g. Shedding vesicles were removed from supernatant by centrifuged at 10,000 g. Finally, supernatant was centrifuged at 110,000g and exosomes were obtained. Isolation of exosomes from serum was performed using ExoQuick Plasma prep and Exosome precipitation kit (SBI, USA).
Tumor formation assay in nude mice NCG (NOD-Prkdc em26Cd52 IL2 rgem26Cd22 /Gpt) mice were purchased from NBRI of Nanjing University (Nanjing, China), which were maintained in a pathogen-free facility. Cells overexpressing SHNG7 were trypsin digested, washed with PBS, and then resuspended in PBS. Then 200µl of the suspended cells (1×10 7 ) were injected into the armpit or peritoneal cavity of each mouse.3-4 weeks later, the mice were sacri ced and tumor weight or metastasis number were examined. All the animal experiments were performed with the approval of the Shandong First Medical University Animal Care and Use Committee.
Statistical Analyses
Each experiment was repeated in triplicate independently. Values were shown as mean ±SD. Statistical analysis was performed using GraphPad Prism 6 (GraphPad Software, USA). P values<0.05 were considered statistically signi cant.
Results
LncRNA SHNG7 is upregulated in HGSOC and correlates with poor outcomes The expression of lncRNA SHNG7 was rstly detected in HGSOC tissues. Compared with normal ovarian epithelium tissues (n=20), lncRNA SHNG7 was signi cantly upregulated in HGSOC tissues (n=30) as detected by qRT-PCR ( Figure 1A). To analyze the effect of lncRNA SHNG7 on the prognosis of HGSOC, the Kaplan-Meier Plotter was applied. The HGSOC patients were divided into two groups, high and low group, based on the median expression value of lncRNA SHNG7. As shown in Figure 1B, the survival rate in high level group was consistently lower than that in low level group.
LncRNA SHNG7 promotes cell proliferation Cell proliferation ability was detected by CCK-8 assay. A2780 cells were transfected with siRNA of lncRNA SHNG7 (or NC) and SKOV3 cells were transfected with lncRNA SHNG7 overexpression vectors (or NC).
The successful transfection was con rmed by qRT-PCR ( Figure 1C). In A2780 cells, the inhibition of lncRNA SHNG7 signi cantly suppressed cell proliferation compared with NC group, and in SKOV3 cells, the overexpression of lncRNA SHNG7 signi cantly promoted cell proliferation compared with NC group ( Figure 1D).
LncRNA SHNG7 promotes cell migration and invasion
To detect the effect of lncRNA SHNG7 on cell migration and invasion, Transwell assay was applied. The Transwell assay results showed that overexpression of lncRNA SHNG7 dramatically promoted cell migration and invasion in SKOV3 cells (Figure 2A). The transfection of siRNA targeting lncRNA SHNG7 dramatically inhibited cell migration and invasion in A2780 cells ( Figure 2B). Then cell cycle distribution of A2780 cells and SKOV3 cells was analyzed using ow cytometry. As shown in Figure 2C, overexpression of lncRNA SNHG7 reduced the proportion of cells in G1 phase and inhibition of lncRNA SNHG5 increased the proportion the G1 phase. In addition, we detected the expression of epithelialmesenchymal transition (EMT)-associated proteins and cell cycle related protein markers (CDK4 and CDK6). The western blot assay showed that overexpression of lncRNA SHNG7 could upregulate the ETM markers such as N-cadherin, β-Catenin, Vimentin, CDK4 and CDK6, and downregulate the epithelial marker E-cadherin ( Figure 2D).
LncRNA SHNG7 could function as ceRNA via directly sponging miR-34a-5p The subcellular location of lncRNA SHNG7 was detected by nuclear mass separation experiment, and we found that lncRNA SHNG7 was mainly located in cytoplasm in A2780 cells ( Figure 3A). According to starBase website, miR-34a-5p was estimated as a potential target of lncRNA SHNG7 ( Figure 3B). To con rm whether lncRNA SHNG7 directly interacted with miR-34a-5p, luciferase reporter assay was performed. As shown in Figure 3C, miR-34a-5p signi cantly reduced the luciferase activity in cells transfected with SHNG7 WT reporter, but not in cells with SHNG7 MT reporter.
LncRNA SHNG7 is carried by exosomes and promotes angiogenesis
Exosomes were isolated from serum and observed by transmission electron microscopy ( Figure 4A). The exosomes was veri ed based on the speci c marker protein (CD63 and TSG101) ( Figure 4B). Then we isolated exosomes from medium of SKOV3 cells transfected with lncRNA SHNG7 overexpression vectors or control vectors. Compared with exosomes from SKOV3 cells (NC), the relative expression of lncRNA SHNG7 in exosomes from SKOV3 cells (SHNG7) was signi cantly increased ( Figure 4C). These results indicated that exosomes are potential carriers of lncRNA SHNG7 in HGSOC. Then we cultured HUVECs with medium from SKOV3 cells (SHNG7) or SKOV3 cells (NC) to detect the role of exosomal lncRNA SHNG7 on angiogenesis in vitro. The ring formation assay showed that HUVECs transfected with exosomal lncRNA SHNG7 had a clear tendency toward ring formation compared with control group ( Figure 4D). And the lncRNA SHNG7 upregulated angiogenesis related protein markers, such as Notch 1, VEGF A and Dll4 ( Figure 4E).
miR-34a-5p directly targets Notch1
Based on the predictive results of TargetScan 7.1, Notch 1 was a putative target gene of miR-34a-5p. One potential miR-34a-5p binding site was found in the 3'UTR of Notch 1 ( Figure 5A). Luciferase reporter was applied to verify the prediction. The luciferase reporter assay showed that miR-34a-5p signi cantly declined the luciferase activity of Notch 1-WT reporter, but no considerable change was observed in Notch 1-MT reporter ( Figure 5B). Immunohistochemical staining was performed to compare the expression level of Notch 1 in the HGSOC tissues (n=30) and normal control samples (n=20). It was found that Notch 1 had higher intensity of immunostaining in HGSOC (76.7%) than in normal control (20.0%) ( Figure 5C and D).
LncRNA SHNG7 enhances the growth and metastasis in vivo Compared with the NC group, the average weight of tumor formed by cells overexpressing SHNG7 was much higher, which indicated that SHNG7 enhanced the growth of ovarian cancer cells in vivo ( Figure 6A and B). Consistently, SHNG7 signi cantly increased the number of peritoneal metastasis in NCG mice ( Figure 6C and D).
Discussion
Recent studies have demonstrated that lncRNAs play an important role in various physiological and pathological processes. In many diseases, especially cancer, lncRNAs are often released from control (13). In the past few years, a large number of lncRNAs have been discovered in mammalian transcriptomes (14). Emerging evidences proved that lncRNA SHNG7 functions as oncogene in human cancers and positively related to clinicopathological characteristics and poor prognosis of patients (15). LncRNAs that are enriched in the cytoplasm could regulate gene expression at the post-transcriptional level via interacting with miRNAs (16). In several cancers, lncRNA SHNG7 was found in the cytoplasm, such as glioblastoma, colorectal cancer (17) and prostate cancer (18). Similarly in this study, we identi ed that lncRNA SHNG7 was located in cytoplasm in A2780 cells. And lncRNA SHNG7 was demonstrated function as a ceRNA to sponge miR-34a-5p and regulate its target gene Notch 1. In HGSOC tissues, lncRNA SHNG7 was upregulated. And the overexpression of lncRNA SHNG7 dramatically promoted cell proliferation, migration and invasion. These evidences indicated that lncRNA SHNG7 act as an oncogene in HGSOC.
Notch signaling pathway is a highly conserved signaling pathway that determines the fate of cells (19). Its receptors and ligands are type I transmembrane proteins that could regulate cell functions through cell-cell interactions (20). In mammals, the Notch signaling pathway includes four receptors (Notch1-Notch4). Notch 1 is an important member of the Notch family (21). Recent studies showed that Notch 1 is not only important for normal cell differentiation (22), its pathophysiological changes are also related to the occurrence and development of some tumors (23). In most cases, activation of Notch signal has oncogenic effects in vitro and animal modules (Capobianco, 1997 #181). In a variety of cancers, Notch 1 is found dysregulated (Nefedova, 2004 #182). Studies have shown that Notch plays an important role in the development of follicles and corpus luteum (Vorontchikhina, 2005 #183). The expression of activated Notch1 and its downstream hes1 gene in ovarian adenocarcinoma is signi cantly higher than that in ovarian adenoma and normal ovarian tissue (Hopfer, 2005 #180), indicating that the Notch signaling pathway (especially Notch 1) is closely related to the development of ovarian cancer.
Between tumor cells or other cells in the body, such as lymphocytes or antigen recognition cells, exosomes rely on the special transmitter components contained in them to transmit special biological signals, thereby affecting the proliferation, migration and invasion ability of tumor cells (Yanez-Mo, 2015 #189). More and more studies have shown that the exosomes in the microenvironment of tumor cells are related to the occurrence and development of tumor cells (Lowry, 2015 #185). These exosomes can promote the proliferation, invasion and metastasis of tumor cells by regulating the body's immune response and promoting neovascularization in tumors, in addition, exosomes secreted by these tumor cells can directly act on other tumor or non-tumor cells (Boelens, 2014 #184). In the present study, we demonstrated the presence of lncRNA SHNG7 in exosomes from serum of HGSOC patients and medium of cancer cells. And the exosomal lncRNA SHNG7 could promote angiopoiesis of HUVECs.
In conclusion, we identi ed the carcinogenic role of lncRNA SHNG7 in HGSOC. The overexpression of lncRNA SHNG7 signi cantly promoted cell proliferation, migration and invasion. LncRNA SHNG7 could be carried to serum or medium from cancer cells, and the exosomal lncRNA SHNG7 dramatically promoted angiopoiesis. Besides, we identi ed that lncRNA SHNG7 could function as a ceRNA against miR-34a-5p, which further regulating Notch 1. The newly identi ed lncRNA SHNG7/ miR-34a-5p / Notch 1 axis provides novel insight into the proliferation and metastasis of HGSOC and represents a potential therapeutic target for the clinical treatment of HGSOC.
Declarations
Ethics approval and consent to participate This study was approved by the Ethics Committee of the First A liated Hospital of Shandong First Medical University. All participants were recruited after prociding a signed informed consent.
Consent for publication
Not available.
Availability of data and materials
The data used and/or analysis during the current study are available from the corresponding author on reasonable request.
Competing interests:
The authors declare that they have no competing interests. Author's contributions This study was conceived, designed and interpreted by ZH. CJ and LX undertook the data acquisition, analysis, and interpretation. YCZ and GWT were responsible for the comprehensive technical support. CJ contributed to the inspection of data and nal manuscript. All authors read and approved the nal manuscript. LncRNA SNHG7 could function as a competing ceRNA for miR-34a-5p. A. LncRNA SNHG7 was mainly located in cytoplasm in A2780 cells. B. The predicted targeting sequence of miR-34a-5p on lncRNA SNHG7. C. Luciferase reporter assay in A2780 cells. LncRNA SNHG5 enhanced tumor growth and metastasis in vivo. A. The photograph of tumors harvested from different groups. B. The tumor weights were compared between two groups (data are mean ± SEM, *p < 0.05, n =10). C. Abdominal metastasis nodes (red arrow) of two group mice. D. SNHG5 increased the number of metastasis (data are mean ± SEM, *p < 0.05, n =3). | 4,118.8 | 2021-11-09T00:00:00.000 | [
"Medicine",
"Biology"
] |
Advances in Radio Science MIMO performance of a planar logarithmically periodic antenna with respect to measured channel matrices
The increasing interest in wireless transmission of highest data rates for multimedia applications (e.g. HDTV) demands the use of communication systems as e.g. described in the IEEE 802.11n draft specification for WLAN including spatial multiplexing or transmit diversity to achieve a constant high data rate and a small outage probability. In a wireless communications system the transmission of parallel data stream leads to multiple input/multiple output (MIMO) systems, whose key parameters heavily depend on the properties of the mobile channel. Assuming an uncorrelated channel matrix the correlation between the multiplexed data streams is caused by the coupling of the antennas, so that the radiation element becomes an even more important part of the system. Previous work in this research area (Klemp and Eul, 2006) has shown that planar log.-per four arm antennas are promising candidates for MIMO applications providing two nearly decorrelated radiators, which cover a wide frequency range including both WLAN bands at 2.4 GHz and 5.4 GHz. Up to now the MIMO performance of this antenna is mainly analyzed by simulations. In this contribution measured channel matrices in a real office environment are studied in terms of the antenna’s MIMO performance such as outage probability. The obtained results recorded by using a commercial platform are compared to the simulated ones.
Introduction
The possibility of achieving a remarkable performance gain of data rate and link reliability by using spatial multiplexing in wireless communication systems lead to the first WLAN standard 802.11n that uses MIMO techniques.However, there are no concrete rules for antenna designers to gener-Correspondence to: H. Rabe (rabe@hft.uni-hannover.de)ate antennas that offer a good MIMO performance preferably in as many channel scenarios as possible.The used antenna structures, the placing as well as the arrangement can have significant influence on the overall system performance.This contribution shows, how polarization diversity can be used to increase orthogonality of the single subchannels from a theoretical point of view by means of a stochastic channel model together with a narrow band assumption (see Sect. 2) for a 2×2 MIMO system.To prove this statement several measurements of channel matrices have been made in an office environment with planar logarithmically periodic antennas (LPs) as well as monopole antennas at typical WLAN frequencies at 2.4 GHz and 5.4 GHz.The measurements have been taken with the HaLo220 testbed that allows the transmission and reception of two RF signals simultaneously.A LP antenna enables the radiation of two linearly polarized waves with an axial ratio of better than 30 dB.The performance differences between a setup with LP antennas and with standard monopole antennas are discussed in Sect.3.2.The paper ends with a conclusion.
MIMO channel
To make reliable MIMO measurements of channel matrices some general boundaries should be taken into account.Sections 2.1 and 2.2 describe the necessity of considering the coherence time and the coherence bandwidth in order to fulfil the narrow band assumption and to preserve stationarity of the channel for short time periods.By considering the aforementioned conditions the channel can be expressed as a channel matrix H consisting of complex values shown in Sect.2.3.This section further discusses the influence of the antennas on the channel matrix separately for LOS (Line Of Sight) and NLOS (Non Line Of Sight) scenarios.
Published by Copernicus Publications on behalf of the URSI Landesausschuss in der Bundesrepublik Deutschland e.V.
H. Rabe et al.: MIMO performance of a planar logarithmically periodic antenna
Coherence time
In general, the impulse response of a radio channel is assumed to be linear but not stationary.The movement of persons and objects in the channel causes more or less rapidly changing wave propagation conditions.In the case of a WLAN system, which is intended to operate in an indoor environment, the velocity of the moving scatterers rarely exceeds v sc =5 km/h (Jiménez, 2002) in most cases.Depending on the center frequency the phase of the transmission path is being changed due to a moving scatterer.Regarding a short time period only, the transmission phase can be assumed to be stationary as well as the whole channel impulse response.The time interval, in which the correlation of the channel response does not rise above 0.5 is defined as coherence time T C in this paper.During the coherence time interval T C the channel impulse response can be described as one complex delta peak that scales the obtained signal in the equivalent baseband domain.As a rule of thumb formula Eq. ( 1) can be used to determine T C .
The coherence time decreases with shrinking wavelength so the worst case is a transmission at the highest used frequency.
For the high WLAN frequency band around 5.4 GHz the coherence time can be estimated to T C =40 ms.However, this formula is not adequate in many cases.Another approach (Rappaport, 1996), that is said to be more appropriate determines a smaller value of T C =18 ms which was assumed to be valid for the examined WLAN channels.
Coherence bandwidth
It is furthermore assumed, that every frequency component in the baseband signal is distorted the same way so the frequency response of the channel is ideally flat.It is important for accurate measurements that this requirement is fulfilled by the radiated signals.Analogous to the coherence time one can determine a coherence bandwidth B C which is directly connected to the delay spread σ τ of the channel.The longer the dominant propagation pathes are the more sensitive the phase difference is between adjacent signal frequencies.The coherence bandwidth can be estimated by expression (2) according to Sklar (2001) and is related to the delay spread σ τ that includes the path lengths of the channel implicitly.
For indoor channels the delay spread strongly depends on the scenario.When transmitter and receiver are both situated in a small room, where the path lengths are supposed to be short, there will be no large spread in the arrival of the signal echoes.Even if the path length are long the delay spread can be short if a strong line of sight component or another dominant propagation path exists.A worst case scenario would be a channel with two similarly weighted pathes of different length, for example a LOS and a reflected component coming from a wall.In this paper the delay spread is assumed to stay below σ τ =100 ns, which is more than sufficient for the office environment where the measurement took place and about the double of the value measured by McDonnel et al. (1998).
Therefore, the coherence bandwidth can be estimated from expression (2) to 200 kHz.This restriction will be considered in Sect.3.
MIMO channel matrix
With the preceding simplifications of a narrow band channel the impulse response h(t) degenerates to one complex value between transmit and receive antenna within a coherence time interval.With multiple transmit and receive antennas each single impulse response value can be composed in the MIMO matrix H as shown in Fig. 1.The receiver signal expands to a receive vector y, which M elements are a linear combination of the N elements in the transmit vector s.The relation between transmit and receive vectors is expressed in Eq. ( 3).
The channel capacity C MIMO in the case of equally divided power P T over the transmit antennas can be expressed as follows according to Foschini and Gans (1998).
E is the identity matrix, λ i is the ith eigenvalue of the matrix product between H and its hermitian H H and σ 2 n is the noise power which is assumed to be white gaussian noise.From Eq. ( 4) one can see that the spectral efficiency is increased by using higher transmit power or by trying to achieve a high rank of the matrix product H H H or the matrix H, respectively.For WLAN indoor channels rayleigh fading is assumed so the matrix H is filled with independent and identically-distributed (IID) complex random variables for each coherence time.Unfortunately it is very hard to achieve the data rate which is offered by the pure rayleigh fading channel model.In fact the signals are transmitted and received at the feeding points of the antennas and will be somehow correlated (see Schumacher and Kermoal, 2002) which avoids reaching the theoretical limit.The model which is used in 802.11n for modeling some of these effects is a stochastic channel model with one or more azimuthal probability density functions describing the incoming and outgoing waves from an antenna array for a few coherence times.
. Rabe et al.: MIMO performance of a planar logarithmically periodic antenna ig. 1.The MIMO channel matrix H itself and the matrix H inluding correlation due to the involved antennas ne or more azimuthal probability density functions describng the incoming and outgoing waves from an antenna array or a few coherence times.Considering the resulting correation of the signals on the feeding point extends the channel atrix H to the matrix H as presented in figure 1.In adition to the stochastic modeling of the WLAN channel the fficial model considers a LOS component which is added to he pure stochastic channel.Therefore the channel matrix H an be thought of being composed from a sum of two matries HLOS and HNLOS with different weights as expressed n equation 5.
he rice factor K determines these weights and is fixed acording to a certain channel scenario while the NLOS comonent is following the fluctuations of the channel modeled y IID values.The power P describes the power of the transitted signals and scales the whole channel matrix.For an ntenna designer it is desirable to benefit from spatial muliplexing offered by the NLOS and LOS component of the hannel as well.The following two sections 2.3.1 and 2.3.2 escribe which antenna parameters are important to look at n order to increase the channel capacity.spectively.In general, the antennas are co izontally placed dipole in the x-y plane w amount of the vertical polarized wave an erty is described by the axial ratio AR, wh the ratio between the maximum and the value of the electric field strengths in the p Assigning a certain AR to the receiving 2 and normalizing the vertically received horizontally received can be determined to this symmetrical case, the same relation c case of receiving a horizontal wave with ented dipole, where the received power is means of the AR.The higher the AR is, pression of the unwanted orthogonal pol is.According to the channel scenario in nel matrix H for the 2x2 case of infinite side and a finite AR on receiver side, the proportional to expression (6).
HLOS ≈ 1 1/AR 1/AR 1
The equation shows that the channel matri and more the identity matrix as the AR inc values of H HH and the spectral efficienc increasing similarly.If an antenna array over a wide angular spread, the channel m high for many LOS cases.Considering the resulting correlation of the signals on the feeding point extends the channel matrix H to the matrix H as presented in Fig. 1.In addition to the stochastic modeling of the WLAN channel the official model considers a LOS component which is added to the pure stochastic channel.Therefore the channel matrix H can be thought of being composed from a sum of two matrices HLOS and HNLOS with different weights as expressed in Eq. ( 5).
HNLOS
(5) The rice factor K determines these weights and is fixed according to a certain channel scenario while the NLOS component is following the fluctuations of the channel modeled by IID values.The power P describes the power of the transmitted signals and scales the whole channel matrix.For an antenna designer it is desirable to benefit from spatial multiplexing offered by the NLOS and LOS component of the channel as well.The following two Sects.2.3.1 and 2.3.2 describe which antenna parameters are important to look at in order to increase the channel capacity.
Antennas in a LOS scenario
The LOS component describes a fixed propagation path between transmit and receive antenna.This fact prevents the use of pattern diversity (see Sect. 2.3.2) in this case, which aims at exciting the eigenpathes of the channel using different propagation pathes.A more promising approach is the use of polarization diversity.Radiating and receiving two orthogonally polarized waves with a low crosstalk allows the transmission of two independent data channels.O performance of a planar logarithmically periodic antenna 3 hannel matrix H itself and the matrix H ine to the involved antennas hal probability density functions describd outgoing waves from an antenna array times.Considering the resulting correon the feeding point extends the channel atrix H as presented in figure 1.In adstic modeling of the WLAN channel the ders a LOS component which is added to hannel.Therefore the channel matrix H eing composed from a sum of two matri-LOS with different weights as expressed etermines these weights and is fixed acchannel scenario while the NLOS comthe fluctuations of the channel modeled ower P describes the power of the transscales the whole channel matrix.For an is desirable to benefit from spatial multhe NLOS and LOS component of the e following two sections 2.3.1 and 2.3.2 nna parameters are important to look at the channel capacity.spectively.In general, the antennas are coupled, so the hor izontally placed dipole in the x-y plane will receive a littl amount of the vertical polarized wave anyway.This prop erty is described by the axial ratio AR, which is expressed a the ratio between the maximum and the minimum absolut value of the electric field strengths in the polarization ellipse Assigning a certain AR to the receiving antennas in figur 2 and normalizing the vertically received power P H to 1 th horizontally received can be determined to P H = 1/AR 2 .I this symmetrical case, the same relation can be shown in th case of receiving a horizontal wave with the vertically ori ented dipole, where the received power is also determined b means of the AR.The higher the AR is, the higher the sup pression of the unwanted orthogonal polarization directio is.According to the channel scenario in figure 2 the chan nel matrix H for the 2x2 case of infinite AR on transmitte side and a finite AR on receiver side, the channel matrix i proportional to expression (6).
The equation shows that the channel matrix approaches mor and more the identity matrix as the AR increases.The eigen values of H HH and the spectral efficiency respectively ar increasing similarly.If an antenna array offers a high AR over a wide angular spread, the channel matrix rank will sta high for many LOS cases.vertical polarized wave with the power P R .The receiving dipole antennas are oriented along the z-and the y-axis, respectively.In general, the antennas are coupled, so the horizontally placed dipole in the x-y plane will receive a little amount of the vertical polarized wave anyway.This property is described by the axial ratio AR, which is expressed as the ratio between the maximum and the minimum absolute value of the electric field strengths in the polarization ellipse.Assigning a certain AR to the receiving antennas in Fig. 2 and normalizing the vertically received power P H to 1 the horizontally received can be determined to P H =1/AR 2 .In this symmetrical case, the same relation can be shown in the case of receiving a horizontal wave with the vertically oriented dipole, where the received power is also determined by means of the AR.The higher the AR is, the higher the suppression of the unwanted orthogonal polarization direction is.According to the channel scenario in Fig. 2 the channel matrix H for the 2×2 case of infinite AR on transmitter side and a finite AR on receiver side, the channel matrix is proportional to expression (6).
HLOS ≈ The equation shows that the channel matrix approaches more and more the identity matrix as the AR increases.The eigenvalues of H H H and the spectral efficiency, respectively, are increasing similarly.If an antenna array offers a high AR over a wide angular spread, the channel matrix rank will stay high for many LOS cases.
Antennas in a NLOS scenario
For the NLOS components of the channel it is necessary to unite the stochastic channel model including angle of arrival (AoA), angle of departure (AoD) and the according angular spreads with the complex far field patterns of the single antenna elements in the array.The measurement of MIMO channel matrices in WLAN channels requires a possibility for simultaneously transmitting and receiving RF signals in both frequency bands around 2.4 GHz and 5.4 GHz.The equipment used and the applied measurement method is explained in section 3.1.The subsection 3.1.1shortly introduces the logarithmic periodic antenna used for the measurement.The measurement results taken in the office environment are described in section 3.1.2are shown and discussed in section 3.2.
Measurement Hardware
The HaLo220 system consists of two router like devices equipped with antennas and connected to a PC via USB as shown in figure 3.Each device contents two dual band transceiver units that can be configured as receiver or transmitter.A fast memory allows the playback or the acquisition of baseband data with different bandwidths up to 40 MHz.The synchronous transmission and reception of data makes it possible to measure the complex impulse response of the channel.However it is necessary to estimate the CIR (channel impulse response) for each path of the 2x2 system by considering the restrictions of the coherence time and bandwidth from section 2. The method applied to separate the single pathes is an evaluation of the transmission of two CW signals with a slight frequency shift ∆f of 200 kHz.The receiver starts recording with a sample rate of f s = 5 MSamples/s, when the signal level crosses an adjustable trigger level as the CW signals impinge.After acquisition, n = 50000 samples of the received complex baseband data are transferred to the PC where an FFT is applied.The duration of the transmission is t tr = 10 ms = n/f s , which lies within the coherence time interval of 18 ms according to section 2.1.The FFT for each of the two receiver signals shows two peaks with different magnitudes and angles according to the transmitted CW signals.By determining these peaks and their values the channel matrix can be composed for every coherence time and the channel capacity can be calculated according to equation (4).Because of the flat fading channel the complex values in the frequency domain result in a delta peak in the envelope correlation coefficient ρ e can be approximated by the square of the absolute value of the correlation coefficient |ρ ij | 2 between the antennas i and j in an antenna array.It determines the degree of linear independency of the columns in the channel matrix and is expressed in formula 7 according to Fujimoto and James.
The variance σ 2 is the part of the radiated or received power which contributes to the transmission in a certain channel scenario.The covariance R ij expresses the amount of power which is commonly radiated or received by the antennas i and j and which can not be distinguished anymore in the receiver signals.The goal is to minimize the covariance R ij which can only be done on antenna side.For the stochastic channel model, the covariance follows Eq. ( 8).
The expression contents the far field patterns C for each antenna i in an array in ϑ and ϕ polarization in dependence of the azimuthal and elevational directions.The patterns of the antenna elements i and j are multiplied for each polarization direction and can be interpreted as the commonly radiated part of the power in one spatial direction.The crosspolarization ratio XP R is a property of the propagation channel and is a measure for the conversion of a linearly polarized wave into its orthogonal polarization.If the ratio is XP R=1, the power of a linearly polarized wave is equally distributed into its co-and cross polarization direction.This is generally assumed in indoor WLAN channels and can be explained by the numerous reflections in the propagation pathes.The product of the patterns is further on multiplied by the propability density function emphasizing the spatial direction of the dominant propagation pathes.Another factor is the exponential term expressing the phase delay between the antenna signals in a specific room direction.As mentioned before it is important to minimize the the covariance between each antenna element in an antenna array for MIMO applications.Three antenna parameters out of Eq. ( 8) can be tuned to achieve this goal.One parameter is the pattern diversity.If the antennas "look" in different directions the product of the patterns is minimized.Another approach is separating the antennas in order to benefit from the phase difference between the signals.The final possibility is to use antennas with different polarizations.In this case, the pattern product can be minimized.If this condition can be kept over a large angular spread, the correlation is small in many channel scenarios.
It was shown, how the LOS and the NLOS components of the channel matrix H can be minimized.In both cases it can be pointed out that polarization diversity helps to exploit the spatial diversity of the channel.This statement has been proven by measured channel matrices for different antenna setups as described in the following section.
Measurement
The measurement of MIMO channel matrices in WLAN channels requires a possibility for simultaneously transmitting and receiving RF signals in both frequency bands around 2.4 GHz and 5.4 GHz.The equipment used and the applied measurement method is explained in Sect.3.1.The Sect. 3.1.1shortly introduces the logarithmic periodic antenna used for the measurement.The measurement results taken in the office environment are described in Sect.3.1.2are shown and discussed in Sect.3.2.
Measurement hardware
The HaLo220 system consists of two router like devices equipped with antennas and connected to a PC via USB as shown in Fig. 3.Each device contents two dual band transceiver units that can be configured as receiver or transmitter.A fast memory allows the playback or the acquisition of baseband data with different bandwidths up to 40 MHz.The synchronous transmission and reception of data makes it possible to measure the complex impulse response of the channel.However it is necessary to estimate the CIR (channel impulse response) for each path of the 2×2 system by considering the restrictions of the coherence time and bandwidth from Sect. 2. The method applied to separate the single pathes is an evaluation of the transmission of two CW signals with a slight frequency shift f of 200 kHz.The receiver starts recording with a sample rate of f s =5 MSamples/s, when the signal level crosses an adjustable trigger level as the CW signals impinge.After acquisition, n=50 000 samples of the received complex baseband data are transferred to the PC where an FFT is applied.The duration of the transmission is t tr =10 ms=n/f s , which lies within the coherence time interval of 18 ms according to Sect.2.1.The FFT for each of the two receiver signals shows two peaks Adv.Radio Sci., 6, [55][56][57][58][59][60][61]2008 www.adv-radio-sci.net/6/55/2008/with different magnitudes and angles according to the transmitted CW signals.By determining these peaks and their values the channel matrix can be composed for every coherence time and the channel capacity can be calculated according to Eq. ( 4).Because of the flat fading channel the complex values in the frequency domain result in a delta peak in the time domain with the same value.To get trustable predictions of the channel behavior, 1000 measurements for each channel scenario have been taken.
Antenna hardware
To demonstrate the influence of polarization diversity on the examined channels two different types of antennas were used.The first type is the well known λ/4 monopole that can be mounted on the HaLo system directly.Two of these elements were used in a distance of about λ for f =2.4 GHz and were both aligned vertically.This setup offers weak polarization diversity.In contrast, the logarithmic periodical antenna shown in Fig. 4 provides the radiation of two linearly polarized waves with high AR.This planar antenna type is a self-complementary structure applied on F R4 substrate consisting of four orthogonally placed arms as shown in Fig. 4 ( Klemp and Eul, 2006).It has originally been designed for broadband applications.Exciting an opposite pair of arms with a differential signal generates a linear polarized wave radiating normally to the substrate plane in both directions.The feeding points for each element are in the center of the antenna and are connected via a semirigid coaxial waveguide from the backside.To reduce interactions with the coaxial cables, the feeding network and the generation of mantle modes a λ/4 reflector is placed on the backside of the antenna.Hence, the gain in main beam direction is increased to approximately G=5 dBi for a narrow band around the center frequency in a full wave simulation at 2.4 GHz.To switch between the two measurement frequencies the distance of the reflector can be adjusted.In this configuration, the simulated AR of the orthogonal arm pairs is better than 30 dB for both frequencies.
Office scenario
A floor plan of the examined office environment is illustrated in Fig. 5.The transmit and receive antennas have been placed in the same height of about 1.2 m above the ground and in a distance of 5 m.To separate between a channel which is dominated by the LOS component and the NLOS component, respectively, a metal plate is placed in front of the transmit antennas.The plate causes a reflection of the transmitted signal to the back of the room.To avoid the direct backscattering from the plate back into the LP antenna it was twisted a little bit (see the left image in Fig. 5).In the monopole setup the LP antennas were just replaced by two vertically oriented monopoles at the same place that are directly connected to the HaLo housing.Three different an-
Antenna Hardware
To demonstrate the influence of polarization diversity the examined channels two different types of antennas we used.The first type is the well known λ/4 monopole th can be mounted on the HaLo system directly.Two of the elements were used in a distance of about λ for f = 2.4 GH and were both aligned vertically.This setup offers weak p larization diversity.In contrast, the logarithmic periodic antenna shown in 3.1.1provides the radiation of two linear polarized waves with high AR.This planar antenna type is self-complementary structure applied on F R4 substrate co sisting of four orthogonally placed arms as shown in figu 3.1.1(Klemp).It has originally been designed for broadba applications.Exciting an opposite pair of arms with a diffe ential signal generates a linear polarized wave radiating no mally to the substrate plane in both directions.The feedi points for each element are in the center of the antenna a are connected via a semirigid coaxial waveguide from t backside.To reduce interactions with the coaxial cables, t feeding network and the generation of mantle modes a λ reflector is placed on the backside of the antenna.Hence, t gain in main beam direction is increased to approximately = 5 dBi for a narrow band around the center frequency in tenna setups were examined.In two cases similar antenna types were used on transmitter and receiver side, monopole and LP antennas respectively.In the third case two different types of antennas were used, an LP antenna on receiver side and two monopoles on transmitter side.
Measurement results
The measurement results are presented as a cumulative distribution function (cdf) of the measured channel capacities for 1000 coherence time intervals.The cdf can be interpreted as the probability of the channel capacity to fall below the value on the abscissa.It is also called the outage rate.Figure 6 shows the results for the LOS scenario (without metal plate) for the two subchannels for each of the three antenna configurations.In the case of using two LP antennas both subchannels are strong.The curves show a high slope because of the dominating static LOS component of the channel.There is no significant change in the channel capacity and due to the good polarization decoupling the channel matrix has high rank which results in a good quality of both subchannels.In comparison to the other antenna setups the sum capacity of the LP ×LP case is the best.Replacing the transmit antenna with two vertically oriented monopoles causes a decrease of the stronger subchannel due to the missing gain into the direction of the receiving antenna.The influence on the weaker subchannel is more significant.It decreases from about 10 bit/s/Hz to 2.5 bit/s/Hz at an outage rate of −1.Both signals are transmitted with vertical polarization and will be received by the vertical polarized arm pair of the LP antenna with more signal power than by the horizontally polarized arm pair.By means of the channel matrix the elements in the second column are much lower than the ones in the first column.The subchannels can hardly be seperated in this case.Most of the signal power received in the horizontal direction is caused by reflections on the walls which avoids the weaker subchannel to vanish completely.This relation results from the smaller slope of the weak compared to the it was twisted a little bit (see the left image in figure 5).In the monopole setup the LP antennas were just replaced by two vertically oriented monopoles at the same place that are directly connected to the HaLo housing.Three different antenna setups were examined.In two cases similar antenna types were used on transmitter and receiver side, monopole and LP antennas respectively.In the third case two different types of antennas were used, an LP antenna on receiver side and two monopoles on transmitter side.
Measurement results
The measurement results are presented as a cumulative distribution function (cdf) of the measured channel capacities for 1000 coherence time intervals.The cdf can be interpreted as the probability of the channel capacity to fall below the value on the abscissa.It is also called the outage rate.Figure 6 shows the results for the LOS scenario (without metal plate) for the two subchannels for each of the three antenna configurations.In the case of using two LP 5).In the monopole setup the LP antennas were just replaced by two vertically oriented monopoles at the same place that are directly connected to the HaLo housing.Three different antenna setups were examined.In two cases similar antenna types were used on transmitter and receiver side, monopole and LP antennas respectively.In the third case two different types of antennas were used, an LP antenna on receiver side and two monopoles on transmitter side.
Measurement results
The measurement results are presented as a cumulative distribution function (cdf) of the measured channel capacities for 1000 coherence time intervals.The cdf can be interpreted as the probability of the channel capacity to fall below the value on the abscissa.It is also called the outage rate.Figure 6 shows the results for the LOS scenario (without metal plate) for the two subchannels for each of the three antenna configurations.In the case of using two LP strong subchannel which implies a higher weight of rayleigh fading in the channel according to a decreasing rice factor K in Eq. ( 5).In the last antenna setup with monopole antennas on transmit and receiver side this effect gets even stronger.
The very constant far field pattern of the single antenna elements in azimuthal direction leads to multiple transmission pathes that are frequently changing.The weaker LOS path results in a SNR loss which can be seen by the vertically shifted curve of the stronger subchannel with and without LP antenna.The higher influence of the stochastic properties of the channel results in a higher spread of the measured channel capacities.Similar results can be obtained by the measurement at 5.4 GHz in Fig. 7.All mentionable effects can be shown regarding the setups of two LPs and two monopoles.In comparison to the measurement at 2.4 GHz the curves are shifted horizontally because of the higher path loss.The shape of the curves, however, is kept as the propagation properties seem to be similar.The subchannels of the LOS dominated channel in the LP ×LP setup is expressed by the high slope of the curves.Both subchannels are strong due to the polarization decoupling which drops in the case of two monopole antennas on each side.Here again the stochastic influences of the channel decrease the slope of the curves.Up to now it can be pointed out that polarization diversity leads to a better decoupling of the subchannels for both frequencies in the LOS case.The good SNR of the channel is due to the gain in main beam direction resulting in a high capacity of both subchannels.shape of the curves, however, is kept as the propagation properties seem to be similar.The subchannels of the LOS dominated channel in the LP x LP setup is expressed by the high slope of the curves.Both subchannels are strong due to the polarization decoupling which drops in the case of two monopole antennas on each side.Here again the stochastic influences of the channel decrease the slope of the curves.Up to now it can be pointed out that polarization diversity leads to a better decoupling of the subchannels for both frequencies in the LOS case.The good SNR of the channel is due to the gain in main beam direction resulting in a high capacity of both subchannels.
The next step is to analyze the results for the NLOS case.According to section 2.3.2 polarization diversity will lead to small correlation of the antenna elements and therefore to a decoupling of the subchannels.The comparison of the subchannel capacities for the setup of two LP antennas and four monopole antennas for 5.4 GHz are shown in figure 8.For both antenna setups the capacities are smaller compared to the ones obtained in the LOS scenarios because of the missing LOS component which has been faded out by placing the metal plate in front of the transmit antennas.Starting with the LP antenna setup it can be determined that the smaller subchannel is not that strong anymore in comparison to one in the LOS case.This is due to the depolarization effect of the channel considered in expression 8 with the variable XP R. The slope of the curves is decreased because of the changing channel conditions as already seen in the changing weight of the Rice factor between the monopole and the LP setup in the LOS case.In the NLOS case the stronger subchannels are similar except for the slope.It seems that in the monopole setup the propagation pathes are more fluctuating than in the LP channel.Although this property would lead to a decor- small correlation of the antenna elements and therefore to a decoupling of the subchannels.The comparison of the subchannel capacities for the setup of two LP antennas and four monopole antennas for 5.4 GHz are shown in Fig. 8.For both antenna setups the capacities are smaller compared to the ones obtained in the LOS scenarios because of the missing LOS component which has been faded out by placing the metal plate in front of the transmit antennas.Starting with the LP antenna setup it can be determined that the smaller subchannel is not that strong anymore in comparison to one in the LOS case.This is due to the depolarization effect of the channel considered in expression (8) with the variable XP R. The slope of the curves is decreased because of the changing channel conditions as already seen in the changing weight of the Rice factor between the monopole and the LP setup in the LOS case.In the NLOS case the stronger subchannels are similar except for the slope.It seems that in the monopole setup the propagation pathes are more fluctuating than in the LP channel.Although this property would lead to a decorrelation in conjunction with the distance of about λ for the monopoles it can be pointed out that the weaker subchannel of the LP setup is still better than in the monpole setup.This emphasizes the role of polarization decoupling in the NLOS case.
Conclusions
The contribution highlighted the antenna properties leading to a good MIMO performance for LOS as well as for NLOS channels assuming a stochastic channel model and a narrow band propagation channel.It has been determined that polarization diversity can improve orthogonality of the subchannels.This issue has been investigated using different antenna setups with high and low polarization decoupling.The LP antenna offers a high depolarization and leads to a good performance compared to a monopole setup that is used in the most present WLAN devices.In future WLAN applications even more antennas will have to be placed in order to exploit the spatial diversity properties of the channel.Due to the unaltered geometrical restrictions antenna designers will be forced to consider every possibility to decrease the correlation between the single antenna elements.Applying planar antennas with different polarization and pattern properties can be a solution of this problem as presented in this paper.
Fig. 1 .
Fig. 1.The MIMO channel matrix H itself and the matrix H including correlation due to the involved antennas.
Figure 2 reveals this relation.A transmit antenna radiates an ideally
Fig. 2 .
Fig. 2. Reception of an ideally linearly polarized wave by a cros dipole with finite axial ratio
Fig. 2 .
Fig. 2. Reception of an ideally linearly polarized wave by a cross dipole with finite axial ratio.
Fig. 5 .
Fig. 5. Floor plan of the office room
Fig. 5 .
Fig. 5. Floor plan of the office room.
Fig. 6 .Fig. 6 .
Fig. 6.Measured cdf of the subchannels at 2.4 GHz for three antenna setups in the LOS scenario antennas both subchannels are strong.The curves show a high slope because of the dominating static LOS component of the channel.There is no significant change in the channel capacity and due to the good polarization decoupling the
Fig. 8 .
Fig. 8. Measured cdf of the subchannels at 5.4 GHz for two antenna setups in the NLOS scenario
Fig. 8 .
Fig. 8. Measured cdf of the subchannels at 5.4 GHz for two antenna setups in the NLOS scenario.
The next step is to analyze the results for the NLOS case.According to Sect.2.3.2 polarization diversity will lead to | 8,824.2 | 2008-05-26T00:00:00.000 | [
"Engineering",
"Physics"
] |
Classroom Teachers’ Awareness, Difficulties and Suggestions about Students with Learning Disabilities in Mathematics
Classroom teachers’ awareness of the characteristics of students with mathematics learning difficulties is important for the planning and implementation of individualized intervention programs for students. This study aims to examine classroom teachers’ awareness of students with mathematics learning difficulties, as well as the mathematics teaching strategies they apply according to their professional knowledge and experience and their approaches towards these children. Case study design, one of the qualitative research methods, was used in the study. The study was conducted with 5 classroom teachers working in 5 different provinces in Turkey in the 2021-2022 academic year. A semi-structured interview form was used to collect the data. The participants of the study were determined by the easily accessible sampling method. The interviews were recorded on a voice recorder with the permission of the classroom teachers, and the data obtained were analyzed by content analysis method. The results of the study show that the knowledge of classroom teachers about the concept of mathematics learning disabilities varies according to the experience and working time of the teachers, the majority of the teachers do not have sufficient knowledge, they do not receive special training about this situation in their undergraduate education and in the institutions where they work, their knowledge about the process of referring students with mathematics learning disabilities to the necessary institutions when they encounter students with mathematics learning disabilities is insufficient, and teachers feel themselves inadequate in the education of these students.
Introduction
Every human being is born with individual differences.Each person learns in line with his/her interest, speed and ability, but it is not yet possible to create an educational environment that prioritizes the individual characteristics of each child in mass education.The education of children who differ significantly from their peers in terms of their individual characteristics is considered within the scope of special education and these children are characterized as individuals with special learning difficulties.In this context, an education program specific to children is planned and implemented (Eripek et al., 1996;Koç & Korkmaz, 2019).
Although the origin of the discussions on the scope and definition of the concept of specific learning disability dates back to the 1930s (İlker & Melekoğlu, 2017), it is defined as the difficulties that arise in the process of acquiring and applying speaking, listening, reading-writing, reasoning and basic mathematical skills (Kirk, 1963;MEB, 2006;Şimşek, 2012).Individuals with specific learning disabilities may experience difficulties in various skills such as mathematical operations, reading, writing, psycho-motor skills, recognizing and combining words, and reading comprehension (Altun et al., 2011).Specific learning disabilities are generally handled under four categories: reading difficulties (dyslexia), mathematics learning difficulties (dyscalculia), written expression difficulties (dysgraphia) and learning disorders that cannot be named (Köroğlu, 2008).
Mathematics Learning Disability (Dyscalculia)
When the literature on mathematics learning disability is examined, it is seen that researchers use different expressions such as mathematical disabilities, arithmetic learning disabilities (Koontz, 1996), mathematics learning disorder or dyscalculia (Morsanyi et al., 2018), and disorder specific to arithmetic skills.Mathematics learning disorder is defined as a deficiency or disorder in various skills such as understanding and seeing numerical and spatial relationships, inadequacy in acquiring mathematical knowledge and skills, understanding mathematical relationships, recognizing and writing symbols, number concept, counting principles and learning arithmetic (Beacham & Trott, 2005;Mutlu, 2016).Köroğlu (2008) states that individuals with learning disabilities in mathematics have difficulties in many areas such as careless, slow and incorrect calculations, difficulty in understanding terms, number symbols and magnitudes, visual perception, time perception, sequencing events and problem steps, recognizing and drawing geometric shapes, understanding fractions, daily life, money and calculations.These children have many common characteristics and common difficulties with other children who have the same problems as them.These common characteristics and difficulties may not be observed at the same rate in all individuals with math learning disabilities.Geary (2011) classified the common characteristics of children with math learning disabilities as shown in the figure below.
Figure 1 Common Characteristics of Individuals with Dyscalculia
Difficulty in understanding numbers: Difficulties in distinguishing the signs of numbers, miscalculation, difficulty in four operation skills/ slow solving, difficulty in understanding and solving problems, difficulty in time perception, difficulty in strategy making skills, difficulty in distinguishing the direction of operations, difficulty in learning fractions.
Difficulty in ordering numbers: Using fingers while doing operations, having difficulty in ordering or comparing numbers (big/small), having difficulty in determining the solution steps of problems, having difficulty in calculating change.
Difficulty in understanding symbols: Deficits in orientation skills, deficits in visual perception (difficulty in recognizing and drawing simple geometric shapes), confusion caused by symbols.
In the literature, there are both national and international studies aiming to determine the level of knowledge of classroom teachers and mathematics teachers about students with dyscalculia and their needs regarding dyscalculia (Saravanabhavan & Saravanabhavan, 2010;Sezer & Akın, 2011;Şimşek & Arslan, 2022;Wadlington & Wadlington, 2006;Wadlington et al., 2006).However, it can be said that there are very few studies on dyscalculia, especially in the national context, and more studies and in-depth information are needed (Baldemir & Tutak, 2022;Sezer & Akın, 2011).Children with mathematics learning disabilities begin to be diagnosed especially in the first years of primary school.At this point, classroom teachers working in primary schools play a vital role in dyscalculia (Başar & Göncü, 2018).In this context, this study is thought to fill an important gap in the literature.
Method
In this study, case study method, one of the qualitative research methods, was used.The most important feature of a case study is to investigate the depth of one or more situations.In this study, classroom teachers' awareness of students with mathematics learning difficulties was tried to be evaluated comprehensively (Yıldırım & Şimşek, 2013).
Participants
In this study, five classroom teachers actively working in the Ministry of National Education were studied.The participants of the study were determined by convenience sampling method.Convenience sampling method is based on items that are available, quick and easy to reach.In this study, the participants were determined as people who volunteered to participate in the study and were easy to reach.The participants were coded as P1, P2, P3, P4, P5 in the order of application (Table 1).Although P1 and P4 had 5 years of experience, teachers with different experiences participated in the study.
Data Collection Tool
In this study, data were collected through semistructured interviews.Three experts were consulted to ensure the content validity of the 8 questions in the interview form.As a result of this interview, it was seen that the questions in the semi-structured interview form were understood by the participants and served the purpose.Interviews were conducted via telephone.The interviews lasted approximately 60 minutes and were recorded with a voice recorder.These records were then transcribed and made ready for analysis.The interview questions are as follows;
Data Analysis
The data analyzed by using the content analysis method."Content analysis is to bring together similar data within the framework of certain themes and to interpret the data by organizing them in a way that the reader can understand" (Yıldırım & Şimşek, 2013, p.259).According to the answers given by the classroom teachers participating in the study, coding was made according to the concepts extracted from the data obtained.Themes were determined in order to collect the codings among certain categories.The data obtained were organized and interpreted according to the emerging themes.
Findings
In this section, classroom teachers' awareness of students with mathematics learning disabilities is presented according to their observations and experiences in the classroom within the framework of interviews with classroom teachers (Figure 2).
Awareness of Teachers and Parents about Special Education
In this section, teachers' definitions of special education, competencies, and family awareness and support are discussed.
Teachers' Definitions of Special Education
When classroom teachers were asked about their thoughts on mathematics learning disability (dyscalculia), all of them stated that they had heard of this concept.While three of the teachers correctly defined dyscalculia, the other two teachers did not have sufficient knowledge.
The teachers stated that dyscalculia is a learning disorder (P2) and that the student has difficulty in reading comprehension, reading, writing and thinking skills (P1).P4 defined dyscalculia as "... a special difference that causes students with normal or above normal intelligence to have low success in reading and writing skills." When the teachers were asked about their views on dyscalculia, it was determined that three of the teachers had heard of this concept before, one teacher could not remember, and the other teacher had never heard of it.The teachers who expressed their opinions about dyscalculia stated that dyscalculia is a learning disability in mathematics (P5, P4), and that this concept is used for students who are less successful than their peers in mathematical numbers, symbols and calculations, in the development of mathematical skills, in problem solving situations, in mathematical reasoning, and who sometimes cannot even use their hands in calculations (P1).
Considering the teacher's experience (P1-5 years) and the frequency of working with special children, it was seen that the teacher's awareness varied.
Teachers' Competencies related to Special Education
Three of the classroom teachers who participated in the study stated that they received training during their undergraduate education, while two of them stated that they did not receive training.The teachers who received training stated that they did not find the training they received sufficient.P5 "...It was mentioned that we should try different ways in learning difficulties and that one way will surely carry that student forward.Although this is basically a correct statement, it would have been better to have a wider range of practical training to guide these paths.In fact, the basis of mathematics is reading comprehension and interpretation.I think students who cannot learn to read with the constructivist method also have problems in mathematics learning."and drew attention to the need for practical training.P2 "... Children should learn with games.There should be separate classes for math.They should learn addition or subtraction with small wooden cylinders, cubes, etc.A garden should be used for calculating the perimeter.By doing and experiencing, all pre-service teachers should first learn by themselves and then educate children."Again, she emphasized the importance of learning applied teaching methods during the training process of pre-service teachers.
When the classroom teachers were asked to what extent they felt adequate for the education of students with mathematics learning difficulties, three teachers felt adequate, while the other two teachers stated that they did not feel adequate.P1 stated that she followed the current resources on the subject and participated in related training seminars, while P4 stated that she did not consider herself sufficient because they did not receive training on how to teach special students in their own classrooms during their undergraduate education.
P3 said, "I consider myself sufficient, but I believe that it would be more efficient if more experts in their field gave one-on-one lessons."P1 said, "I follow current sources, I try to attend training seminars on this subject, I try to read publications on dyslexia or dyscalculia."P5 said, "There will be times when I am not sufficient; but I see those times as learning opportunities for myself by struggling." Although teacher support and dedication is a very important factor in the education of children with dyscalculia, it is not sufficient alone.These children should receive special education support from special education institutions outside the classroom in line with their needs, and classroom teachers should support the student in the light of this educational planning and support the progress of the process.
Family Awareness and Support
According to the opinions of classroom teachers about the responsibilities of the families of students with dyscalculia, the importance of families in the development of students' academic achievement is great.They stated that these students with dyscalculia need more attention and special education support, and in cases where the family does not provide the necessary support, more burden is placed on the shoulders of the teacher.A significant number of classroom teachers drew attention to the cooperation between family, school and teacher.They stated that they observed that the child's academic achievement increased with the cooperation provided.However, they stated that if the necessary cooperation is not provided by the family, the child is deprived of the special education he/she should receive.
In cases where there is no family support, teachers stated that families are indifferent, unconscious or do not accept their child's problem for various reasons.They stated that if the family knows about the problem but is indifferent to this situation, they take care of the child themselves.P5 suggested that families should be given training on the subject in order to raise their awareness.P1 stated that some families unfortunately tend to hide or ignore their children's special situations in order to prevent their children from being exposed to some stigmatization.P3 stated that some families realize the situation and direct their children to other fields, but some families do not accept the situation and think that their children are lazy.According to the teachers' views, considering that individuals with dyscalculia need special education and support, their academic development is negatively affected when they lack family support, which is the most important pillar of this support.
Difficulties Experienced by Students in Need of Special Education
All of the classroom teachers participating in the study stated that they had students with specific learning difficulties in their classrooms.Classroom teachers responded to the question "What are the characteristics of students with mathematics learning disabilities?"by expressing the difficulties experienced by students with dyscalculia in mathematics.When Figure 3 is examined, the difficulties experienced by the students are discussed in two categories: cognitive and affective difficulties.Regarding the low level of intelligence, the teachers said, "Children with mathematics learning difficulties have less developed areas of mathematical intelligence or they learn later."(P3) and "they have difficulty in moving from abstract to concrete" (P5).Mathematical difficulties were discussed in two categories: difficulties with mathematical symbols, operations and concepts and difficulties with mathematical skills.Under the title of difficulties related to mathematical symbols, operations and concepts, it was observed that the classroom teachers emphasized the following characteristics: not knowing the meaning of symbols, not being able to match symbols with the number of objects, difficulty in learning 4 operations skills and multiplication tables, difficulty in rhythmic counting, writing numbers backwards, not being able to recognize geometric shapes, difficulty in reading and writing natural numbers, difficulty in perceiving time and difficulty in understanding fractions.For example, P2 said about a student, "He does not know the concept of number.He cannot learn numbers.Numbers are just shapes for him.In other words, the child does not know that the number two corresponds to two fingers.He can write numbers up to 20, but he does not know what these numbers correspond to.When I ask him to show 5 fingers, he can show up to three or four, but he cannot comprehend that 5 fingers correspond to the number 5." The teacher stated that the student could not match the symbols with the number of objects.Teachers also emphasized the lack of mathematical skills under the category of mathematical difficulties.They stated that students had difficulties in problem solving, critical thinking, reasoning and associating mathematics with daily life, other disciplines and concepts.For example, P1 stated that students "lag behind their peers in numerical reasoning, logic and reasoning skills, and critical thinking compared to their peers."and P5 stated that "Students cannot forget and use what they have learned when they move on to another subject, so they cannot establish relationships between concepts."They expressed as follows.P1 stated, "In their daily lives, for example, when they go to the market, they have trouble calculating change and calculating how much money they have spent.These problems also cause social phobia in children."P1 emphasized that students have difficulty in associating mathematics with daily life.
Affective Challenges
In addition to cognitive difficulties, classroom teachers reported that students experienced affective difficulties.The findings show that students face many affective difficulties such as distraction, hyperactivity, anxiety, lack of self-confidence, introversion, negative attitude towards mathematics, fear of mathematics and learned helplessness.Classroom teachers stated that children with learning disabilities are generally withdrawn, experience loss of self-confidence due to exposure to negative behaviors, and have problems expressing themselves.They stated that they show learned helplessness and negative attitudes towards mathematics, especially because their failures in mathematics are perceived negatively by their environment.Some of the teachers stated that these children experience two extremes; some children are very quiet and do not want to participate in class, while others are overly active and show negative behaviors towards their friends.The reason for this situation is the negative behaviors that the child is exposed to by his/her family or teacher and the constant feeling of failure.P1 stated that "I observe a phobia and helplessness towards mathematics, a decrease in interest in other subjects because of not being able to realize oneself in mathematics, as well as "distraction and hyperactivity".
Recommendations for the Education of Students in Need of Special Education
The classroom teachers participating in the study were asked about their suggestions for the education of students in need of special education.The findings were categorized under two main headings: suggested teaching approach and suggested methods and techniques.
In the context of teaching approach, the classroom teachers stated that the education of students with special needs should be done directly by special education specialists, the needs of the child should be determined, positive discrimination should be made for these students, individual time should be allocated in and out of class, students should be given additional time, self-confidence should be gained with a positive approach and cooperation with the family should be made.P2 stated that "special time should be allocated for these children and special education specialists should provide education" and P4 stated that students "should receive education according to the individualized education program".P1 stated that "These children should be treated sensitively and positive discrimination should be made in some issues.In a class of 40 students, if you explain to these children the way you explain to other children, you will lose the child in the class.You would be ignoring him/her" and stated that positive discrimination should be given to students.
In the context of suggested methods and techniques, they stated that students can be taught outside the classroom, students can be taught especially with the support of concrete materials, students should repeat frequently, approaches such as gamification, visualization, one-to-one education and learning by doing and experiencing can be used in the education to be given to students.For example, P2 said, "You should not be able to teach something without drawing shapes, without using materials, without involving the child.These children should receive special education.Apart from the education in the classroom, a certain hour should be allocated to them and frequent repetitions should be made with materials by making mathematics completely concrete."P2 emphasized the importance of concrete material-supported education.P5, on the other hand, stated the following about the need for frequent repetition: "When the subject is given, plenty of examples should be given, frequent repetition should be done and enough time should be given to the student".Another issue emphasized by the teachers is the planning of education according to the needs of the child.Teacher P5 made suggestions such as determining the points where students are deficient in detail and creating an individual plan, creating environments where they can feel success, paying attention to the transition from simple to complex, from near to far, from concrete to abstract since they are slow in high-level thinking skills, and associating them with life by giving more examples.
P1, P2 and P5 teachers stated that the child who was subjected to the right educational interventions with regular and planned cooperation reached the level of his/her peers, and that the child's mental state improved after a certain period of time by using educational methods appropriate to the child's needs in a systematic way.P1 said, "I witnessed that the child who was subjected to the right educational interventions with regular and planned cooperation reached the level of his/her peers.I saw that after being diagnosed with specific learning disabilities and receiving special education for a few years, he/she did not receive the same diagnosis in the institution where he/she went for re-diagnosis.For example, I had a 3rd grade student who had difficulty learning the multiplication table.I was observing behavioral disorders along with high math anxiety in the child.We started to work on multiplication tables with the gamification method.For example, while playing hopscotch, I would include the multiplication table.When the child reached the 5th grade, I observed that he memorized the multiplication table and his success in other courses increased significantly", and drew attention to the fact that the academic performance of students increased with the right educational intervention.
Conclusion, Discussion and Suggestions
In this study, classroom teachers' awareness of students with mathematics learning difficulties, the mathematics teaching strategies they apply according to their professional knowledge and experience, and their approaches towards these children were examined.As a result of the study, although there were some teachers who had never heard of the concept of dyscalculia, the majority of the teachers were able to define this concept, albeit partially.Similarly, Büyükkarcı and Akgün-Giray (2023) conducted a study with prospective classroom teachers and found that prospective teachers knew the concept of dyscalculia, although not in depth.Some of the teachers confused the concepts of specific learning disabilities and mathematics learning disabilities (dyscalculia).As a result of the interviews with the classroom teachers who participated in the study, the majority of the teachers who expressed opinions about the concept of dyscalculia expressed this concept as students' inability to perform basic four operations, inability to understand/learn mathematical problems, inability to comprehend numbers and digits, and inability to understand abstract concepts.Similarly in the literature, learning disabilities in mathematics are classified as difficulties in number perception, accurate and fluent calculation, reasoning and problem solving (American Psychiatric Association (APA), 2013).
The opinions of classroom teachers about the causes of dyscalculia are intelligence level, prejudice against mathematics, distraction and hyperactivity, self-confidence problems, and inability to understand commands.Although students with dyscalculia do not have any disadvantage in terms of intelligence level, they are children who do not find the motivation, interest and experiences necessary for teaching, and the reasons for the special learning disability they experience are independent of intelligence (Görgün & Melekoğlu, 2019).Teachers who state that students' mathematical intelligence does not develop have misinformation about individuals with dyscalculia.Attention deficit and hyperactivity, self-confidence problems, behavioral disorders, problems in social relations can be seen together with learning disabilities, but they do not constitute a learning disability on their own.Learning disabilities can be seen together with some emotional and mental reasons, but they are not a direct result of these reasons (National Joint Committee on Learning Disabilities Definition of Learning Disabilities (NJCLD), 1990).As it is understood from this definition, it is understood that the knowledge of classroom teachers about the characteristics of individuals with dyscalculia is not sufficient.Considering that early diagnosis is an important factor in dyscalculia, it is thought that this situation may negatively affect the diagnosis process of students.Classroom teachers should know how to approach a student who has any of the specific learning difficulties in their class.For this reason, having information about both the individual and general characteristics of individuals with dyscalculia will help the planning and implementation processes of individual education programs to be prepared for students to be successful.
In the light of the opinions on the responsibilities of the families of individuals with dyscalculia, classroom teachers stated that students with dyscalculia should be supported by their families and that family, school and teacher cooperation is important.The majority of classroom teachers stated that although family support is important, most families do not provide support.According to the opinions, most of the families think that the teacher should take this responsibility.Classroom teachers attributed the reason for this situation to families' lack of interest, lack of awareness or low awareness.In addition, even if some families are aware of the special situation of their students, they insist on not accepting the situation and not providing special education in order not to be criticized.Teachers should also be aware of their responsibilities in order to ensure cooperation between family, school and teacher.The classroom teacher has a great role in informing the family about dyscalculia, trying to establish cooperation with the family and supporting the family in education.
When the opinions of classroom teachers about the methods and techniques they use in teaching mathematics to students with dyscalculia were examined, some of the classroom teachers stated that they use visual elements and make abstract concepts concrete with materials.Some of the teachers stated that they work one-on-one with their students with dyscalculia.Individualization of teaching in terms of the academic development of students with learning disabilities yields positive results.Considering that individuals with learning disabilities need different areas of education, supporting these individuals sufficiently can prevent students from isolating and comparing themselves with their peers.Similarly, in the study conducted by Altun and Uzuner (2016), according to the opinions of classroom teachers on the education of students with specific learning difficulties, teachers stated that they are interested in such students one-on-one and try to do activities according to the needs of the student.In addition, teachers stated that they try to show a positive approach and interest to these students, repeat the subject, avoid complex expressions and get help from the guidance service when necessary.
Classroom teachers listed the situations in which students with mathematics learning difficulties have difficulty in learning mathematics as writing numbers and numbers backwards, four operations skills, understanding mathematical symbols, recognizing geometric shapes, multiplication table, reading and writing natural numbers and problems.
In the studies conducted on the subject, individuals with dyscalculia have difficulty in mathematical calculations that require complex operations.With the increase in the complexity level of mathematics questions, the inadequacy in students' memory and learning strategies causes negative effects on their learning performance (Bender, 2014).Students with dyscalculia have difficulty in following the steps used in math problems.In order to improve students' ability to solve mathematical problems, long problems with many steps should be divided into shorter and more meaningful steps, important parts of the question should be underlined with colored pencils, and shapes should be used in problem solving (Sezer & Akın, 2011).Other subjects in which students have difficulty are additionsubtraction, multiplication and division.Teachers stated that students had difficulty in multiplication operations as the number of digits increased, they could not memorize the multiplication table, they had difficulty in rhythmic counting, they could not establish the addition-multiplication relationship, and they had difficulties in addition operations with the product.However, they did not mention the number line and fractions, which is one of the most common problems experienced by students.
Classroom teachers listed their observations about the causes of mathematics anxiety in students with mathematics learning difficulties as teacher and parental reactions, parental expectations, lack of self-confidence, comparison with peers, long computation time and feeling of failure.Mathematics anxiety, which is one of the biggest obstacles to learning mathematics, causes individuals to show low study performance and does not allow them to reveal their current potential.When we evaluate mathematics anxiety in terms of children with dyscalculia, it is seen that these children have very high anxiety due to their inability to do mathematics.These children have normal and above normal intelligence levels.
In some cases, their dyscalculia causes them to be exposed to unconscious attitudes and behaviors by both parents and teachers.With the right teaching approaches, the anxiety levels of children with learning difficulties should be reduced, and the child should be able to use mathematical skills supported by teaching strategies specific to this field.
The teaching approaches used by classroom teachers to reduce mathematics anxiety caused by the above-mentioned reasons for students with mathematics learning difficulties are listed as subthemes in the form of fun activities, gaining selfconfidence, giving extra time and having a positive approach.The teaching methods exhibited by the teachers support the literature studies.Implementation of mathematics curricula with daily life activities, fun and educational animations, integrating the lesson with games, etc. will help reduce the anxiety of children with dyscalculia (Geist, 2010).
The majority of classroom teachers stated they did not consider themselves adequate for the education of students with dyscalculia.The reason for this is that they did not receive a special education specific to the field in undergraduate education, each child has different characteristics and they are constantly learning new things from them.Classroom teachers stated that they could not pay much attention to special students because of the crowded classrooms and the necessity of raising the curriculum, and that the education programs were far above the level of these children.Participant teachers stated that students with dyscalculia should receive support from special education specialists, guidance services and family members.Similarly, Bevan and Butterworth (2002) reported that classroom teachers found mathematics curricula to be difficult and complex and therefore inappropriate for students with dyscalculia.
Two of the teachers stated that children could make progress at their own level if they were taught with the right teaching strategies, but they could not reach the level of their peers.Other teachers, on the other hand, stated that children who are subjected to the right educational interventions with regular and planned cooperation reach the level of their peers, and that the child's mental state improves after a certain period of time by using educational methods appropriate to the child's needs in a systematic manner.Being aware of the subjects in which students with dyscalculia have difficulties, knowing the characteristics of individuals with dyscalculia, having knowledge about how to approach these students, and implementing individually tailored education plans will help students progress in the area in which they have difficulty and increase their performance (Mutlu & Aygün, 2020).
Based on the findings obtained from the research, the following suggestions can be made; • Courses related to special learning disabilities and dyscalculia, one of its subclasses, should be given in undergraduate education.It is thought that the number of students with dyscalculia is too high to be ignored.• In-service trainings can be given to classroom teachers to increase their awareness about specific learning disabilities and math learning disabilities and to improve their professional knowledge.• Trainings can be given to families about special learning disabilities and dyscalculia in order to ensure school, family and teacher cooperation. | 7,012.6 | 2023-10-01T00:00:00.000 | [
"Mathematics",
"Education"
] |
Angular analysis of $B^0_d \rightarrow K^{*}\mu^+\mu^-$ decays in $pp$ collisions at $\sqrt{s}= 8$ TeV with the ATLAS detector
An angular analysis of the decay $B^0_d \rightarrow K^{*}\mu^+\mu^-$ is presented, based on proton-proton collision data recorded by the ATLAS experiment at the LHC. The study is using 20.3 fb$^{-1}$ of integrated luminosity collected during 2012 at centre-of-mass energy of $\sqrt{s}=8$ TeV. Measurements of the $K^{*}$ longitudinal polarisation fraction and a set of angular parameters obtained for this decay are presented. The results are compatible with the Standard Model predictions.
Introduction
Flavour-changing neutral currents (FCNC) have played a significant role in the construction of the Standard Model of particle physics (SM). These processes are forbidden at tree level and can proceed only via loops, hence are rare. An important set of FCNC processes involve the transition of a b-quark to an sµ + µ − final state mediated by electroweak box and penguin diagrams. If heavy new particles exist, they may contribute to FCNC decay amplitudes, affecting the measurement of observables related to the decay under study. Hence FCNC processes allow searches for contributions from sources of physics beyond the SM (hereafter referred to as new physics). This analysis focuses on the decay B 0 d → K * 0 (892)µ + µ − , where K * 0 (892) → K + π − . Hereafter, the K * 0 (892) is referred to as K * and charge conjugation is implied throughout, unless stated otherwise. In addition to angular observables such as the forward-backward asymmetry A FB 1, there is considerable interest in measurements of the charge asymmetry, differential branching fraction, isospin asymmetry, and ratio of rates of decay into dimuon and dielectron final states, all as a function of the invariant mass squared of the dilepton system q 2 . All of these observable sets can be sensitive to different types of new physics that allow for FCNCs at tree or loop level. The BaBar, Belle, CDF, CMS, and LHCb collaborations have published the results of studies of the angular distributions for B 0 d → K * µ + µ − [1][2][3][4][5][6][7][8]. The LHCb Collaboration has reported a potential hint, at the level of 3.4 standard deviations, of a deviation from SM calculations [3,4] in this decay mode when using a parameterization of the angular distribution designed to minimise uncertainties from hadronic form factors. Measurements using this approach were also reported by the Belle Collaboration [8] and they are consistent with the LHCb experiment's results and with the SM calculations. This paper presents results following the methodology outlined in Ref. [3] and the convention adopted by the LHCb Collaboration for the definition of angular observables described in Ref. [9]. The results obtained here are compared with theoretical predictions that use the form factors computed in Ref. [10].
This article presents the results of an angular analysis of the decay B 0 d → K * µ + µ − with the ATLAS detector, using 20.3 fb −1 of pp collision data at a centre-of-mass energy √ s = 8 TeV delivered by the Large Hadron Collider (LHC) [11] during 2012. Results are presented in six different bins of q 2 in the range 0.04 to 6.0 GeV 2 , where three of these bins overlap. Backgrounds, including a radiative tail from B 0 d → K * J/ψ events, increase for q 2 above 6.0 GeV 2 , and for this reason, data above this value are not studied.
The operator product expansion used to describe the decay B 0 d → K * µ + µ − encodes short-distance contributions in terms of Wilson coefficients and long-distance contributions in terms of operators [12]. Global fits for Wilson coefficients have been performed using measurements of B 0 d → K * µ + µ − and other rare processes. Such studies aim to connect deviations from the SM predictions in several processes to identify a consistent pattern hinting at the structure of a potential underlying new-physics Lagrangian, see Refs. [13][14][15]. The parameters presented in this article can be used as inputs to these global fits.
Analysis method
Three angular variables describing the decay are defined according to convention described by the LHCb Collaboration in Ref. [9]: the angle between the K + and the direction opposite to the B 0 d in the K * centre-of-mass frame (θ K ); the angle between the µ + and the direction opposite to the B 0 d in the dimuon centre-of-mass frame (θ L ); and the angle between the two decay planes formed by the K π and the dimuon systems in the B 0 d rest frame (φ). For B 0 d mesons the definitions are given with respect to the negatively charged particles. Figure 1 illustrates the angles used. The angular differential decay rate for B 0 d → K * µ + µ − is a function of q 2 , cos θ K , cos θ L and φ, and can be written in several ways [16]. The form to express the differential decay amplitude as a function of the angular parameters uses coefficients that may be represented by the helicity or transversity amplitudes [17] and is written as2 1 dΓ/dq 2 d 4 Γ d cos θ L d cos θ K dφdq 2 = 9 32π 3(1 − F L ) 4 sin 2 θ K + F L cos 2 θ K + 1 − F L 4 sin 2 θ K cos 2θ L −F L cos 2 θ K cos 2θ L + S 3 sin 2 θ K sin 2 θ L cos 2φ +S 4 sin 2θ K sin 2θ L cos φ + S 5 sin 2θ K sin θ L cos φ +S 6 sin 2 θ K cos θ L + S 7 sin 2θ K sin θ L sin φ +S 8 sin 2θ K sin 2θ L sin φ + S 9 sin 2 θ K sin 2 θ L sin 2φ .
Here F L is the fraction of longitudinally polarised K * mesons and the S i are angular coefficients. These angular parameters are functions of the real and imaginary parts of the transversity amplitudes of B 0 d decays into K * µ + µ − . The forward-backward asymmetry is given by A FB = 3S 6 /4. The predictions for the S parameters depend on hadronic form factors which have significant uncertainties at leading order. It is possible to reduce the theoretical uncertainty in these predictions by transforming the S i using ratios constructed to cancel form factor uncertainties at leading order. These ratios are given by Refs. [17,18] as All of the parameters introduced, F L , S i and P ( ) j , may vary with q 2 and the data are analysed in q 2 bins to obtain an average value for a given parameter in that bin.
The ATLAS detector, data, and Monte Carlo samples
The ATLAS experiment at the LHC is a general-purpose detector with a cylindrical geometry and nearly 4π coverage in solid angle [19]. It consists of an inner detector (ID) for tracking, a calorimeter system and a muon spectrometer (MS). The ID consists of silicon pixel and strip detectors, with a straw-tube transition radiation tracker providing additional information for tracks passing through the central region of the detector.3 The ID has a coverage of |η| < 2.5, and is immersed in a 2T axial magnetic field generated by a superconducting solenoid. The calorimeter system, consisting of liquid argon and scintillator-tile sampling calorimeter subsystems, surrounds the ID. The outermost part of the detector is the MS, which employs several detector technologies in order to provide muon identification and a muon trigger. A toroidal magnet system is embedded in the MS. The ID, calorimeter system and MS have full azimuthal coverage.
The data analysed here were recorded in 2012 during Run 1 of the LHC. The centre-of-mass energy of the pp system was √ s = 8 TeV. After applying data-quality criteria, the data sample analysed corresponds to an integrated luminosity of 20.3 fb −1 . A number of Monte Carlo (MC) signal and background event samples were generated, with b-hadron production in pp collisions simulated with P 8.186 [20,21]. The AU2 set of tuned parameters [22] is used together with the CTEQ6L1 PDF set [23]. The EvtGen 1.2.0 program [24] is used for the properties of b-and c-hadron decays. The simulation included modelling of multiple interactions per pp bunch crossing in the LHC with P soft QCD processes. The simulated events were then passed through the full ATLAS detector simulation program based on G 4 [25,26] and reconstructed in the same way as data. The samples of MC generated events are described further in Section 5. 3 ATLAS uses a right-handed coordinate system with its origin at the nominal interaction point (IP) in the centre of the detector and the z-axis along the beam pipe. The x-axis points from the IP to the centre of the LHC ring, and the y-axis points upward. Cylindrical coordinates (r, Φ) are used in the transverse plane, Φ being the azimuthal angle around the z-axis. The pseudorapidity is defined in terms of the polar angle θ as η = − ln tan(θ/2).
Event selection
Several trigger signatures constructed from the MS and ID inputs are selected based on availability during the data-taking period, prescale factor and efficiency for signal identification. Data are combined from 19 trigger chains where 21%, 89% or 5% of selected events pass at least one trigger with one, two, or at least three muons identified online in the MS, respectively. Of the events passing the requirement of at least two muons, the largest contribution comes from the chain requiring one muon with a transverse momentum p T > 4 GeV and the other muon with p T > 6 GeV. This combination of triggers ensures that the analysis remains sensitive to events down to the kinematic threshold of q 2 = 4m 2 µ , where m µ is the muon mass. The effective average trigger efficiency for selected signal events is about 29%, determined from signal MC simulation.
Muon track candidates are formed offline by combining information from both the ID and MS [27]. Tracks are required to satisfy |η| < 2.5. Candidate muon (kaon and pion) tracks in the ID are required to satisfy p T > 3.5 (0.5) GeV. Pairs of oppositely charged muons are required to originate from a common vertex with a fit quality χ 2 /NDF < 10.
Candidate K * mesons are formed using pairs of oppositely charged kaon and pion candidates reconstructed from hits in the ID. Candidates are required to satisfy p T (K * ) > 3.0 GeV. As the ATLAS detector does not have a dedicated charged-particle identification system, candidates are reconstructed with both possible Kπ mass hypotheses. The selection implicitly relies on the kinematics of the reconstructed K * meson to determine which of the two tracks corresponds to the kaon. If both candidates in an event satisfy selection criteria, they are retained and one of them is selected in the next step following a procedure described below. The Kπ invariant mass is required to lie in a window of twice the natural width around the nominal mass of 896 MeV, i.e. in the range [846, 946] MeV. The charge of the kaon candidate is used to assign the flavour of the reconstructed B 0 d candidate.
The B 0 d candidates are reconstructed from a K * candidate and a pair of oppositely charged muons. The fourtrack vertex is fitted and required to satisfy χ 2 /NDF < 2 to suppress background. A significant amount of combinatorial, B 0 d , B + , B 0 s and Λ b background contamination remains at this stage. Combinatorial background is suppressed by requiring a B 0 d candidate lifetime significance τ/σ τ > 12.5, where the decay time uncertainty σ τ is calculated from the covariance matrices associated with the four-track vertex fit and with the primary vertex fit. Background from final states partially reconstructed as B → µ + µ − X accumulates at invariant mass below the B 0 d mass and contributes to the signal region. It is suppressed by imposing an asymmetric mass cut around the nominal B 0 d mass, 5150 MeV < m K πµµ < 5700 MeV. The high-mass sideband is retained, as the parameter values for the combinatorial background shapes are extracted from the fit to data described in Section 5. To further suppress background, it is required that the angle Θ, defined between the vector from the primary vertex to the B 0 d candidate decay vertex and the B 0 d candidate momentum, satisfies cos Θ > 0.999. Resolution effects on cos θ K , cos θ L and φ were found to have a negligible effect on the ATLAS B 0 s → J/ψφ analysis [28]. It is assumed to also be the case for On average 12% of selected events in the data have more than one reconstructed B 0 d candidate. The fraction is 17% for signal MC samples and 2-10% for exclusive background MC samples. A two-step selection process is used for such events. For 4% of these events it is possible to select a candidate with the smallest value of the B 0 d vertex χ 2 /NDF. However, the majority, about 96%, of multiple candidates arise from four-track combinations where the kaon and pion assignments are ambiguous. As these candidates have degenerate values for the B 0 d candidate vertex χ 2 /NDF, a second selection step is required. The B 0 d candidate reconstructed with the smallest value of |m K π − m K * |/σ(m K π ) is retained for analysis, where m K π is the K * candidate mass, σ(m K π ) is the uncertainty in this quantity, and m K * is the world average value of the K * mass.
The selection procedure results in an incorrect flavour tag (mistag) for some signal events. The mistag probability of a B 0 d (B 0 d ) meson is denoted by ω (ω) and is determined from MC simulated events to be 0.1088 ± 0.0005 (0.1086 ± 0.0005). The mistag probability varies slightly with q 2 such that the difference ω − ω remains consistent with zero. Hence the average mistag rate ω in a given q 2 bin is used to account for this effect. If a candidate is mistagged, the values of cos θ L , cos θ K and φ change sign, while the latter two are also slightly shaped by the swapped hadron track mass hypothesis. Sign changes in these angles affect the overall sign of the terms multiplied by the coefficients S 5 , S 6 , S 8 and S 9 (similarly for the corresponding P ( ) parameters) in Equation (1). The corollary is that mistagged events result in a dilution factor of (1 − 2 ω ) for the affected coefficients.
The region q 2 ∈ [0.98, 1.1] GeV 2 is vetoed to remove any potential contamination from the φ(1020) resonance. The remaining data with q 2 ∈ [0.04, 6.0] GeV 2 are analysed in order to extract the signal parameters of interest. Two K * cc control regions are defined for B 0 d decays into K * J/ψ and K * ψ(2S), respectively as q 2 ∈ [8,11] and [12,15] GeV 2 . The control samples are used to extract values for nuisance parameters describing the signal probability density function (pdf) from data as discussed in Section 5.3.
For q 2 < 6 GeV 2 the selected data sample consists of 787 events and is composed of signal B 0 d → K * µ + µ − decay events as well as background that is dominated by a combinatorial component that does not peak in m K πµµ and does not exhibit a resonant structure in q 2 . Other background contributions are considered when estimating systematic uncertainties. Above 6 GeV 2 the background contribution increases significantly, including events coming from B 0
Maximum-likelihood fit
Extended unbinned maximum-likelihood fits of the angular distributions of the signal decay are performed on the data for each q 2 bin. The discriminating variables used in the fit are m K πµµ , the cosines of the helicity angles (cos θ K and cos θ L ), and φ. The likelihood L for a given q 2 bin is where N is the total number of events, the sum runs over signal and background components, n l is the fitted yield for the l th component, n is the sum over n l , and P kl is the pdf evaluated for event k and component l.
In the nominal fit, l iterates only over one signal and one background component. The p are parameters of interest (F L , S i ) and θ are nuisance parameters. The remainder of this section discusses the signal model (Section 5.1), treatment of background (Section 5.2), use of K * cc decay control samples (Section 5.3), fitting procedure and validation (Section 5.4).
Signal model
The signal mass distribution is modelled by a Gaussian distribution with the width given by the per-event uncertainty in the K π µµ mass, σ(m K πµµ ), as estimated from the track fit, multiplied by a unit-less scale factor ξ, i.e. the width given by ξ · σ(m K πµµ ). The mean values of the B 0 d candidate mass (m 0 ) and ξ of the signal Gaussian pdf are determined from fits to data in the control regions as described in Section 5.3. The simultaneous extraction of all coefficients using the full angular distribution of Equation (1) requires a certain minimum signal yield and signal purity to avoid a pathological fit behaviour. A significant fraction of fits to ensembles of simulated pseudo-experiments do not converge using the full distribution. This is mitigated using trigonometric transformations to fold certain angular distributions and thereby simplify Equation (1) such that only three parameters are extracted in one fit: F L , S 3 and one of the other S parameters. Following Ref. [3], the transformations listed below are used: F L , S 3 , S 8 , P 8 : On applying transformation (7), (8), (9), and (10), the angular variable ranges become respectively. A consequence of using the folding schemes is that S 6 (A FB ) and S 9 cannot be extracted from the data. For these schemes the angular parameters of interest, denoted by p in Equation (6), are (F L , S 3 , S i ) where i = 4, 5, 7, 8. These translate into (F L , P 1 , P j ), where j = 4, 5, 6, 8, using Equation (5).
The values and uncertainties of F L and S 3 obtained from the four fits are consistent with each other and the results reported are those found to have the smallest systematic uncertainty.
Three MC samples are used to study the signal reconstruction and acceptance. Two of them follow the SM prediction for the decay angle distributions taken from Ref.
[29], with separate samples generated for B 0 d and B 0 d decays. The third MC sample has F L = 1/3 and the angular distributions are generated uniformly in cos θ K , cos θ L and φ. The samples are used to study the effect of potential mistagging and reconstruction differences between particle and antiparticle decays and for determination of the acceptance. The acceptance function is defined as the ratio of reconstructed and generated distributions of cos θ K , cos θ L , φ, i.e. it is compensating for the bias in the angular distributions resulting from triggering, reconstruction and selection of events. It is described by sixth-order (second-order) polynomial distributions for cos θ K and cos θ L (φ) and is assumed to factorise for each angular distribution, i.e. using ε(cos θ K , cos θ L , φ) = ε(cos θ K )ε(cos θ L )ε(φ). A systematic uncertainty is assessed in order to account for this assumption. The acceptance function multiplies the angular distribution in the fit, i.e. the signal pdf is where g(cos θ K , cos θ L , φ) is an angular differential decay rate resulting from one of the four folding schemes applied to Equation (1) and G(m K πµµ ) is the signal mass distribution. The MC sample generated with uniform cos θ K , cos θ L and φ distributions is used to determine the nominal acceptance functions for each of the transformed variables defined in Equations (7)- (10). The other samples are used to estimate the related systematic uncertainty. Among the angular variables the cos θ L distribution is the most affected by the acceptance. This is a result of the minimum transverse momentum requirements on the muons in the trigger and the larger inefficiency to reconstruct low-momentum muons, such that large values of | cos θ L | are inaccessible at low q 2 . As q 2 increases, the acceptance effects become less severe. The cos θ K distribution is affected by the ability to reconstruct the K π system, but that effect shows no significant variation with q 2 . There is no significant acceptance effect for φ. Figure 2 shows the acceptance functions used for cos θ K and cos θ L for two different q 2 ranges for the nominal angular distribution given in Equation (1).
Background modes
The fit to data includes a combinatorial background component that does not peak in the m K πµµ distribution. It is assumed that the background pdf factorises into a product of one-dimensional terms. The mass distribution of this component is described by an exponential function and second-order Chebychev polynomials are used to model the cos θ K , cos θ L and φ distributions. The values of the nuisance parameters describing these shapes are obtained from fits to the data independently for each q 2 bin.
Inclusive samples of bb → µ + µ − X and cc → µ + µ − X decays and eleven exclusive B 0 d , B 0 s , B + and Λ b background samples are studied in order to identify contributions of interest to be included in the fit model, or to be considered when estimating systematic uncertainties. The relevant exclusive modes found to be of interest are discussed below. Events with B c decays are suppressed by excluding the q 2 range containing the J/ψ and ψ(2S), and by charm meson vetoes discussed in Section 7. The exclusive background decays considered for the signal mode are These background contributions are accounted for as systematic uncertainties estimated as described in Section 7.
Two distinct background contributions not considered above are observed in the cos θ K and cos θ L distributions. They are not accounted for in the nominal fit to data, and are treated as systematic effects. A peak is found in the cos θ K distribution near 1.0 and appears to have contributions from at least two distinct sources. One of these arises from misreconstructed B + decays, such as B + → K + µµ and B + → π + µµ. These decays can be reconstructed as signal if another track is combined with the hadron to form a K * candidate in such a way that the event passes the reconstruction and selection. The second contribution comes from combinations of two charged tracks that pass the selection and are reconstructed as a K * candidate. These fake K * candidates accumulate around cos θ K of 1.0 and are observed in the Kπ mass sidebands away from the K * meson. They are distinct from the structure of expected S-, Pand D-wave Kπ decays resulting from a signal B 0 d → K π µµ transition. The origin of this source of background is not fully understood. The observed excess may arise from a statistical fluctuation, an unknown background process, or a combination of both. Systematic uncertainties are assigned to evaluate the effect of these two background contributions, as described in Section 7.
Another peak is found in the cos θ L distribution near values of ±0.7. It is associated with partially reconstructed B decays into final states with a charm meson. This is studied using Monte Carlo simulated events for the decays D 0 → K − π + , D + → K − π + π + and D + s → K + K − π + . Events with a B meson decaying via an intermediate charm meson D 0 , D + or D + s are found to pass the selection and are reconstructed in such a way that they accumulate around 0.7 in | cos θ L |. These are removed from the data sample when estimating systematic uncertainties, as described in Section 7.
K * c c control sample fits
The mass distribution obtained from the simulated samples for K * cc decays and the signal mode, in different bins of q 2 , are found to be consistent with each other. Values of m 0 and ξ for B 0 d → K * J/ψ and B 0 d → K * ψ(2S) events are used for the signal pdf and extracted from fits to the data. An extended unbinned maximum-likelihood fit is performed in the two K * cc control region samples. There are three exclusive backgrounds included: Λ b → Λcc, B + → K + cc and B 0 s → K * cc. The K * cc pdf has the same form as the signal model, combinatorial background is described by an exponential distribution, and double and triple Gaussian pdfs determined from MC simulated events are used to describe the exclusive background contributions. A systematic uncertainty is evaluated by allowing for 0, 1, 2 and 3 exclusive background components. The control sample fit projections for the variant of the fit including all three exclusive backgrounds can be found in Figure 3.
Fitting procedure and validation
A two-step fit process is performed for the different signal bins in q 2 . The first step is a fit to the K π µ + µ − invariant mass distribution, using the event-by-event uncertainty in the reconstructed mass as a conditional variable. For this fit, the parameters m 0 and ξ are fixed to the values obtained from fits to data control samples as described in Section 5.3. A second step adds the (transformed) cos θ K , cos θ L and φ variables to the likelihood in order to extract F L and the S parameters along with the values for the nuisance parameters related to the combinatorial background shapes. Some nuisance parameters, namely m 0 , ξ, signal and background yields, and the exponential shape parameter for the background mass pdf, are fixed to the results obtained from the first step.
The fit procedure is validated using ensembles of simulated pseudo-experiments generated with the F L and S parameters corresponding to those obtained from the data. The purpose of these experiments is to measure the intrinsic fit bias resulting from the likelihood estimator used to extract signal parameters. These ensembles are also used to check that the uncertainties extracted from the fit are consistent with expectations. Ensembles of simulated pseudo-experiments are performed in which signal MC events are injected into samples of background events generated from the likelihood. The signal yield determined from the first step in the fit process is found to be unbiased. The angular parameters extracted from the nominal fits have biases with magnitudes ranging between 0.01 and 0.04, depending on the fit variation and q 2 bin. A similar procedure is used to estimate the effect of neglecting S-wave contamination in the data sample. Neglecting the S-wave component in the fit model results in a bias between 0.00 and 0.02 in the angular parameters. Similarly, neglecting exclusive background contributions from Λ b , B + and B 0 s decays that peak in m K πµµ near the B 0 d mass results in a bias of less than 0.01 on the angular parameters. All these effects are included in the systematic uncertainties described in section 7. The P ( ) parameters are obtained using the fit results and covariance matrices from the second fit along with Equations (2)-(5).
Results
The event yields obtained from the fits are summarised in Table 1 where only statistical uncertainties are reported. Figures 4 through 9 show for the different q 2 bins the distributions of the variables used in the fit for the S 5 folding scheme (corresponding to the transformation of Equation (8)) with the total, signal and background fitted pdfs superimposed. Similar sets of distributions are obtained for the three other folding schemes: S 4 , S 7 and S 8 . The results of the angular fits to the data in terms of the S i and P ( ) j can be found in Tables 2 and 3. Statistical and systematic uncertainties are quoted in the tables. The distributions of F L and the S i parameters as a function of q 2 are shown in Figure 10 and those for P ( ) j are shown in Figure 11. The correlations between F L and the S i parameters and between F L and the P ( ) j are given in Appendix A.
Systematic uncertainties
Systematic uncertainties in the parameter values obtained from the angular analysis come from several sources. The methods for determining these uncertainties are based either on a comparison of nominal and modified fit results, or on observed fit biases in modified pseudo-experiments. The systematic uncertainties are symmetrised. The most significant ones are described in the following, in decreasing order of importance.
• A systematic uncertainty is assigned for the combinatorial K π (fake K * ) background peaking at cos θ K values around 1.0 obtained by comparing results of the nominal fit to that where data above cos θ K = 0.9 are excluded from the fit.
• A systematic uncertainty is derived to account for background arising from partially reconstructed B → D 0 /D + /D + s X decays, that manifest in an accumulation of events at | cos θ L | values around 0.7. Two-track or three-track combinations are formed from the signal candidate tracks, and are reconstructed assuming the pion or kaon mass hypothesis. A veto is then applied for events in which a track combination has a mass in a window of 30 MeV around the D 0 , D + or D + s meson mass. Similarly, a veto is implemented to reject B + → K + µ + µ − and B + → π + µ + µ − events that pass the event selection. Here B + candidates are reconstructed from one of the hadrons from the K * candidate and the muons in the signal candidate. Signal candidates that have a three-track mass within 50 MeV of the B + mass are excluded from the fit. A few percent of signal events are removed when applying these vetoes, with a corresponding effect on the acceptance distributions. The fit results obtained from the data samples with vetoes applied are compared to those obtained from the nominal fit and the change in each result is taken as the systematic uncertainty from these backgrounds. This systematic uncertainty dominates the measurement of F L at higher values of q 2 .
• The background pdf shape has an uncertainty arising from the choice of model. For the mass distribution it is assumed that an exponential function model is adequate; however, for the angular variables the data are re-fitted using third-order Chebychev polynomials. The change from the nominal result is taken as the uncertainty from this source.
• The acceptance function is assumed to factorise into three separate components, one each for cos θ K , cos θ L and φ. To validate this assumption, the signal MC data are fitted with the acceptance function obtained from that sample. Differences in the fit results from expectation are small and taken as the uncertainty resulting from this assumption.
• A systematic uncertainty is assigned for the angular background pdf model by comparing the nominal result to that with a reduced fit range of m K πµµ ∈ [5200, 5700] MeV, in particular to account for possible residues of the partially reconstructed B-decays.
• A correction is applied to the data by shifting the track p T according to the uncertainties arising from biases in rapidity and momentum scale. The change in results obtained is ascribed to the uncertainty in the ID alignment and knowledge of the magnetic field.
• The maximum-likelihood estimator used is intrinsically biased. Ensembles of MC simulated events are used in order to ascertain the bias in the extracted values of the parameters of interest. The bias is assigned as a systematic uncertainty.
• The p T spectrum of B 0 d candidates observed in data is not accurately reproduced by the MC simulation. This difference in the kinematics results in a slight modification of the acceptance functions. This is accounted for by reweighting signal MC simulated events to resemble the p T spectrum found in data. The change in fitted parameter values obtained due to the reweighting is taken as the systematic uncertainty resulting from this difference.
• The signal decay mode is resonant K * → Kπ decay, but scalar contributions from non-resonant Kπ transitions may also exist. The LHCb Collaboration reported an S-wave contribution at the level of 5% of the signal [4,30]. Ensembles of MC simulated events are fitted with 5% of the signal being drawn from an S-wave sample of events and the remaining 95% from signal. The observed change in fit bias is assigned as the systematic uncertainty from this source. Any variation in S-wave content as a function of q 2 would not significantly affect the results reported here.
• The values of the nuisance parameters of the fit model obtained from MC control samples and fits to the data mass distribution have associated uncertainties. These parameters include m 0 , ξ, the signal and background yields, the shape parameter of the combinatorial background mass distribution, and the parameters of the signal acceptance functions. The uncertainty in the value of each of these parameters is varied independently in order to assess the effect on parameters of interest. This source of uncertainty has a small effect on the measurements reported here.
• Background from exclusive modes peaking in m K πµµ is neglected in the nominal fit. This may affect the fitted results and is accounted for by computing the fit bias obtained when embedding MC simulated samples of Λ b → Λ(1520)µ + µ − , Λ b → pK − µ + µ − , B + → K ( * )+ µ + µ − and B 0 s → φµ + µ − into ensembles of pseudo-data generated from the fit model containing only combinatorial background and signal components. The change in fit bias observed when adding exclusive backgrounds is taken as the systematic error arising from neglecting those modes in the fit.
• The difference from nominal results obtained when fitting the B 0 d signal MC events with the acceptance function for B 0 d is taken as an upper limit of the systematic error resulting from event migration due to mistagging the B 0 d flavour.
• The parameters S 5 and S 8 , as well as the respective P ( ) j parameters are affected by dilution and thus have a multiplicative scaling applied to them. This dilution factor depends on the kinematics of the K * decay and has a systematic uncertainty associated with it. The effect of data/MC differences in the p T spectrum of B 0 d candidates on the mistag probability was studied and found to be negligible. The uncertainty due to the limited number of MC events is used to compute the statistical uncertainty of ω and ω. Studies of MC simulated events indicate that there is no significant difference between the mistag probability for B 0 d and B 0 d events and the analysis assumes that the average mistag probability provides an adequate description of this effect. The magnitude of the mistag probability difference, |ω − ω|, is included as a systematic uncertainty resulting from this assumption.
The total systematic uncertainties of the fitted S i and P ( ) j parameter values are presented in Tables 2 and 3, where the dominant contributions for F L come from the modelling of the angular distributions of the combinatorial background and the partially reconstructed decays peaking in cos θ K and cos θ L . These contributions and in addition also ID alignment and magnetic field calibration affect S 3 (P 1 ). The largest systematic uncertainty contribution to S 3 (P 1 ) comes from partially reconstructed decays entering the signal region. This also affects the measurement of S 5 (P 5 ) and S 7 (P 6 ). The partially reconstructed decays peaking in cos θ L affect the measurement of S 4 (P 4 ) and S 8 (P 8 ), whereas the fake K * background in cos θ K affects S 4 (P 4 ), S 5 (P 5 ), and S 8 (P 8 ). The parameterization of the signal acceptance is another significant systematic uncertainty source for S 4 (P 4 ). The systematic uncertainties are smaller than the statistical uncertainties for all parameters measured.
Comparison with theoretical computations
The Figure 10 for the S parameters, and in Figure 11 for the P ( ) parameters, along with the results presented here.4 QCD factorisation is used by DHMV and JC, where the latter focus on the impact of long-distance corrections using a helicity amplitude approach. The CFFMPSV group takes a different approach, using the QCD factorisation framework to perform compatibility checks of the LHCb data with theoretical predictions. This approach also allows information from a given experimentally measured parameter of interest to be excluded in order to make a fit-based prediction of the expected value of that parameter from the rest of the data.
With the exception of the P 4 and P 5 measurements in q 2 ∈ [4.0, 6.0] GeV 2 and P 8 in q 2 ∈ [2.0, 4.0] GeV 2 there is good agreement between theory and measurement. The P 4 and P 5 parameters have statistical correlation of 0.37 in the q 2 ∈ [4.0, 6.0] GeV 2 bin. The observed deviation from the SM prediction of P 4 and P 5 is for both parameters approximately 2.7 standard deviations (local) away from the calculation of DHMV for this bin. The deviations are less significant for the other calculation and the fit approach. All measurements are found to be within three standard deviations of the range covered by the different predictions. Hence, including experimental and theoretical uncertainties, the measurements presented here are found to agree with the predicted SM contributions to this decay. Figure 11: The measured values of P 1 , P 4 , P 5 , P 6 , P 8 compared with predictions from the theoretical calculations discussed in the text (Section 8). Statistical and total uncertainties are shown for the data, i.e. the inner mark indicates the statistical uncertainty and the total error bar the total uncertainty.
Conclusion
The results of an angular analysis of the rare decay B 0 d → K * µ + µ − are presented. This flavour-changing neutral current process is sensitive to potential new-physics contributions. The B 0 d → K * µ + µ − analysis presented here uses a total of 20.3 fb −1 of √ s = 8 TeV pp collision data collected by the ATLAS experiment at the LHC in 2012. An extended unbinned maximum-likelihood fit of the angular distribution of the signal decay is performed in order to extract the parameters F L , S i and P ( ) j in six bins of q 2 . Three of these bins overlap in order to report results in ranges compatible with other experiments and phenomenology studies. All measurements are found to be within three standard deviation of the range covered by the different predictions. The results are also compatible with the results of the LHCb, CMS and Belle collaborations.
Appendix A Correlation Matrices
Four folding schemes are applied to the data in order to extract F L , S 3 , S 4 , S 5 , S 7 and S 8 from four separate fits. The P ( ) parameters are subsequently derived from the fit results using Equations (2)-(5). It is not possible to extract a full correlation matrix between fitted parameters obtained from different fits. In order to reconstruct the correlation matrix, ensembles of pseudo-experiments are simulated using the pdf corresponding to the nominal angular distributions. Each simulated ensemble has the four folding schemes applied to it and four fits are performed on the resulting samples. The distributions obtained for pairs of parameters obtained from fits to these ensembles are used to compute Pearson correlation coefficients for those pairs. Correlation matrices for F L and the S parameters are reconstructed from all possible pairings for a given q 2 bin. A similar method is used to extract the correlation matrices for the P ( ) parameters. This procedure is repeated for each q 2 bin studied in order to obtain correlation matrices given in the remainder of this appendix. The correlation matrices are statistical only. Contributions from systematic uncertainties are not included, since the measurement precision is statistically limited. | 9,846.2 | 2018-05-10T00:00:00.000 | [
"Physics"
] |
Using Green, Economical, Efficient Two-Dimensional (2D) Talc Nanosheets as Lubricant Additives under Harsh Conditions
Two-dimensional (2D) nanomaterials have attracted much attention for lubrication enhancement of grease. It is difficult to disperse nanosheets in viscous grease and the lubrication performances of grease under harsh conditions urgently need to be improved. In this study, the 2D talc nanosheets are modified by a silane coupling agent with the assistance of high-energy ball milling, which can stably disperse in grease. The thickness and size of the talc nanosheet are about 20 nm and 2 µm. The silane coupling agent is successfully grafted on the surface of talc. Using the modified-talc nanosheet, the coefficient of friction and wear depth can be reduced by 40% and 66% under high temperature (150 °C) and high load (3.5 GPa), respectively. The enhancement of the lubrication and anti-wear performance is attributed to the boundary adsorbed tribofilm of talc achieving a repairing effect of the friction interfaces, the repairing effect of talc on the friction interfaces. This work provides green, economical guidance for developing natural lubricant additives and has great potential in sustainable lubrication.
Introduction
With the development of advanced machines, the working conditions of machines become more and more complicated, due to extreme pressure, high temperature, irradiation [1][2][3][4][5][6][7], etc. Friction and wear are the main reasons for the machine failure. Using lubricants is one of the most widely used strategies for friction and wear reduction [8][9][10][11]. Hence, it is of great importance to develop high-efficiency lubricants. Grease is always used in harsh working conditions, but the performance of grease is greatly affected by temperature and load [12]. At high temperatures, grease is usually oxidized to generate unwanted compounds, which damage the lubrication performance, or even leads to the failure of lubrication [13].
Although lubricating additives takes only a small proportion of grease, the tribological properties of grease depend on the additives to a great extent [13]. Hence, it is of great importance to develop novel additives with excellent performance for greases under harsh conditions. In recent years, nanomaterials as lubricant additives have drawn attention from a large number of researchers, due to their excellent antifriction and anti-wear performance [14,15]. Nanoparticles such as Cu, Fe 3 O 4 , and TiO 2 can effectively reduce friction and wear, especially under boundary lubrication regimes [14][15][16][17]. Two-dimensional (2D) nanomaterials, due to their excellent self-lubricating properties, have attracted much attention [18][19][20]. The typical two-dimensional (2D) nano additives are MoS 2 , WS 2 , graphene, talc, etc. These 2D materials, with one or a few atomic layers, have excellent tribological properties [21][22][23][24]. Because of the high surface energy, nanomaterials usually show a strong Nanomaterials 2022, 12, 1666 2 of 13 tendency to aggregate in lubricants [1,25,26]. In addition, the additive will be recognized as ecofriendly additives for anti-friction and/or anti-wear in lubricating systems if the additive will not release SAPS (sulfated ash, phosphorous and sulfur; SAPS would cause air pollution such as acid rain and haze weather) [27][28][29][30]. Thus, green nanomaterials have been considered as the promising alternative to the typical additive of zinc dialkyl dithiophosphates (ZDDPs), which contain phosphorus and sulfur.
Talc is a 2D layered and naturally abundant mineral; thus, it is low cost, ecofriendly, and stable [31,32]. Its specific gravity ranges from 2.7 to 2.8, and it offers high chemical inertness. Talc layers are weakly bonded with each other through Van der-Waals forces forming the lamellar structure, which facilitates easy shearing, giving it a good self-lubrication performance because the lamellar structure facilitates easy shearing [32][33][34]. Therefore, 2D talc nanosheets are considered as one of the best candidate materials for the development of green lubricating additives. Talc is widely used as solid lubricant with the properties of high crystallinity, low electrical conductivity, high thermal stability and good adsorption properties. It has been found that using talc as lubricant additives can enhance the lubrication of commercial engine oil. Talc also shows potential as an eco-friendly extreme pressure lubricating additive [31,32]. However, it is difficult to disperse talc nanomaterials in viscous grease, and their enhancement for the lubrication of grease under harsh conditions is unknown.
In this study, to improve the dispersion of talc in grease, 2D talc nanosheets are prepared via high-energy ball milling with the assistance of a silane coupling agent. The modified talc can be stably dispersed in grease and has much better tribological properties than those of non-modified talc. The micro structures are characterized and the lubrication performances are studied under harsh conditions. The lubrication mechanism has also been discussed. This work provides green, economical guidance for developing natural lubricant additives and has great potential in sustainable lubrication.
Raw Materials and Instruments
The commercial talc, noted as unmodified talc was purchased from Shanghai Yuanjiang Chemical Co., Ltd. (Shanghai, China). Lithium-based lubricating grease is number 0 of Shangbo, noted as base grease was from Sinopec Lubricating Oil Co., Ltd. (Beijing, China). The silane coupling agent (KH550) came from Nanjing Daoning Chemical Co., Ltd. (Nanjing, China). All materials were the analysis reagents. The ball crusher for preparing the samples was from Changsha Miqi Technology Manufacturing Co, Ltd. (Changsha, China). The friction tests were conducted on the SRV-4 tester (Optimal Instruments, Schwabisch Hall, Germany). The friction pairs were in a ball-disc contact form. The material of both the upper and lower parts of the friction pair was bearing steel (GCr15) with a roughness of 20 nm and a hardness of 650-700 HV. To study the influence of temperature on the tribological properties of talc-grease grease, the temperature was varied from 80 • C to 175 • C under 150 N. To study the influence of load on the tribological properties of talcgrease grease, the load was varied from 150 N to 550 N under 80 • C. The experiments were repeated three times for each condition. The friction conditions were summarized in Table 1. The morphologies of the nanosheets were characterized by scanning electron microscopy (SEM; FEI Quanta 200 FEG, Eindhoven, The Netherlands) and high-resolution transmission electron microscopy (TEM; JEM-2010, Tokyo, Japan) with an accelerating voltage of 120 kV. The surface elements were determined by energy dispersive spectroscopy (EDS; Oxford X-MaxN, Oxford, UK) with an accelerating voltage of 15 kV. The crystal lattice structure and the orientation were identified by X-ray diffraction (XRD) with a Bruker D8 Advance diffractometer (Bruker, Billerica, MA, USA). The chemical structure was obtained by a Fourier Transform Infrared Spectroscopy (FTIR, Vertex 70 V, NETZSCH, Selb, Germany). Three-dimensional (3D) White-light interferometry microscopy (Nexview, ZYGO Lamda, Middleton, CT, USA) was used to detect the widths and depths of wear tracks.
Preparation of Grease
Mechanochemical method from ball milling is by virtue of mechanical energy to induce chemical reactions for preparing new materials or modifying materials [35][36][37]. In this experiment, a silane coupling agent was selected as the modifier of talc powder during ball milling processing. As shown in Figure 1, talc powder and the silane coupling agent were weighed as 1:1, and placed in the ball-grinding tank. The mass ratio of grinding the ball (Φ 4 mm) to talc nanosheets was 50:1. The speed of the ball mill was 400 r/min. The mill time was 4 h for the full reaction between the talc and the silane coupling agent. After cooling to room temperature, the samples were washed three times with acetone centrifugation. The centrifuge was rotated at 5000 r/min for 5 min. Filtering was then performed after the centrifugation. The modified talc was put in a vacuum chamber at 50 • C for 1 h. After that, the modified talc was added to the grease. The mass ratio of milling ball to the talc/grease mixture was 5:1. The speed of the ball mill process had a speed of 200 r/min for 4 h. After filtering, the prepared, modified talc-based grease was achieved. Via the same ball mill process without the agent, the unmodified talc-based grease was also obtained. copy (Nexview, ZYGO Lamda, Middleton, CT, USA) was used to detect the widths and depths of wear tracks.
Preparation of Grease
Mechanochemical method from ball milling is by virtue of mechanical energy to induce chemical reactions for preparing new materials or modifying materials [35][36][37]. In this experiment, a silane coupling agent was selected as the modifier of talc powder during ball milling processing. As shown in Figure 1, talc powder and the silane coupling agent were weighed as 1:1, and placed in the ball-grinding tank. The mass ratio of grinding the ball (Ф 4 mm) to talc nanosheets was 50:1. The speed of the ball mill was 400 r/min. The mill time was 4 h for the full reaction between the talc and the silane coupling agent. After cooling to room temperature, the samples were washed three times with acetone centrifugation. The centrifuge was rotated at 5000 r/min for 5 min. Filtering was then performed after the centrifugation. The modified talc was put in a vacuum chamber at 50 °C for 1 h. After that, the modified talc was added to the grease. The mass ratio of milling ball to the talc/grease mixture was 5:1. The speed of the ball mill process had a speed of 200 r/min for 4 h. After filtering, the prepared, modified talc-based grease was achieved. Via the same ball mill process without the agent, the unmodified talc-based grease was also obtained.
Characterization of Talc
The morphological characteristics of the unmodified talc and modified talc are well characterized by SEM and TEM. As shown in Figure 2, the unmodified talc exhibits an obvious 2D layered structure, which is uniformly distributed within the range of about 2-3 µm. There are some wrinkled and exfoliated structures obtained for the talc via ball milling (Figure 2b). From the TEM image (Figure 2c,d), the talc shows relatively integrated lamellar structure and the layer thickness is about 20 nm. According to the EDS analysis ( Figure 2e,f), the modified talc with the silane coupling agent has a higher fraction of Si than that of unmodified talc, which means that the agent has successfully grafted to the talc surface.
Characterization of Talc
The morphological characteristics of the unmodified talc and modified talc are well characterized by SEM and TEM. As shown in Figure 2, the unmodified talc exhibits an obvious 2D layered structure, which is uniformly distributed within the range of about 2-3 μm. There are some wrinkled and exfoliated structures obtained for the talc via ball milling (Figure 2b). From the TEM image (Figure 2c,d), the talc shows relatively integrated lamellar structure and the layer thickness is about 20 nm. According to the EDS analysis (Figure 2e,f), the modified talc with the silane coupling agent has a higher fraction of Si than that of unmodified talc, which means that the agent has successfully grafted to the talc surface. Figure 3 shows the XRD pattern of the talc powder. It can be seen that the characteristic diffraction peaks of the modified talc and the unmodified talc are consistent with that of the standard talc (Mg 3 Si 4 O 10 (OH) 2 ) spectrum, i.e., 9.45 [38]. Thus, the modified talc keeps an ordered crystal structure. The undamaged 2D structure of the talc with easy shearing effect will make a great contribution to lubrication performance.
Nanomaterials 2022, 12, x FOR PEER REVIEW 5 of 13 Figure 3 shows the XRD pattern of the talc powder. It can be seen that the characteristic diffraction peaks of the modified talc and the unmodified talc are consistent with that of the standard talc (Mg3Si4O10(OH)2) spectrum, i.e., 9.45°, 19.32°, 28.59°,36.21°, 60.50° and 59.937° [38]. Thus, the modified talc keeps an ordered crystal structure. The undamaged 2D structure of the talc with easy shearing effect will make a great contribution to lubrication performance. In addition, the chemical structure of the modified talc is analyzed by FTIR (Figure 4). The modified talc shows the absorption peak at 667 cm −1 , which is resulted from MgO (Mg=O stretching) in talc [39]. The obvious absorption peak of Si-O-Si stretching at 1017 cm −1 means the talc surface is successfully modified by the agent via a covalent bond [40,41]. Thus, the dispersion and lubrication performance of the nanosheets in grease have a great potential to be realized. In addition, the chemical structure of the modified talc is analyzed by FTIR (Figure 4). The modified talc shows the absorption peak at 667 cm −1 , which is resulted from MgO (Mg=O stretching) in talc [39]. The obvious absorption peak of Si-O-Si stretching at 1017 cm −1 means the talc surface is successfully modified by the agent via a covalent bond [40,41]. Thus, the dispersion and lubrication performance of the nanosheets in grease have a great potential to be realized.
The results of TG curves of talc and grease are shown in Figure 5. The temperature ranges from 30 • C to 700 • C with a heating rate of 10 • C/min in nitrogen environment. Compared with the unmodified talc and the agent, it can be concluded that the weight loss of the modified talc (5%) within the temperature range between 30 and 180 • C results from the release of the grafted agents. The weight loss of unmodified talc and modified talc are only 2.4% and 8.9%, respectively, until about 500 • C, indicating the talc nanosheet is very thermostable. This is one reason that the talc improves the thermal stability of grease as shown in Figure 5b. The results of TG curves of talc and grease are shown in Figure 5. The temperature ranges from 30 °C to 700 °C with a heating rate of 10 °C/min in nitrogen environment. Compared with the unmodified talc and the agent, it can be concluded that the weight loss of the modified talc (5%) within the temperature range between 30 and 180 °C results from the release of the grafted agents. The weight loss of unmodified talc and modified talc are only 2.4% and 8.9%, respectively, until about 500 °C , indicating the talc nanosheet is very thermostable. This is one reason that the talc improves the thermal stability of grease as shown in Figure 5b.
Lubrication Performance
Friction from the moving machine pairs generates a lot of heat and the local instantaneous temperature can reach as high as 300 °C [14]. The high-temperature tribological tests are more representative of the actual operating conditions in industrial machinery. So it is crucial to study the effect of high temperature on lubrication properties of grease. It can be seen from Figure 6a-c that the COF of modified talc-based grease (0.5 wt.%) is The results of TG curves of talc and grease are shown in Figure 5. The temperature ranges from 30 °C to 700 °C with a heating rate of 10 °C/min in nitrogen environment. Compared with the unmodified talc and the agent, it can be concluded that the weight loss of the modified talc (5%) within the temperature range between 30 and 180 °C results from the release of the grafted agents. The weight loss of unmodified talc and modified talc are only 2.4% and 8.9%, respectively, until about 500 °C , indicating the talc nanosheet is very thermostable. This is one reason that the talc improves the thermal stability of grease as shown in Figure 5b.
Lubrication Performance
Friction from the moving machine pairs generates a lot of heat and the local instantaneous temperature can reach as high as 300 °C [14]. The high-temperature tribological tests are more representative of the actual operating conditions in industrial machinery. So it is crucial to study the effect of high temperature on lubrication properties of grease. It can be seen from Figure 6a-c that the COF of modified talc-based grease (0.5 wt.%) is
Lubrication Performance
Friction from the moving machine pairs generates a lot of heat and the local instantaneous temperature can reach as high as 300 • C [14]. The high-temperature tribological tests are more representative of the actual operating conditions in industrial machinery. So it is crucial to study the effect of high temperature on lubrication properties of grease. It can be seen from Figure 6a-c that the COF of modified talc-based grease (0.5 wt.%) is always stable and lower at high temperature compared with the one without talc. Due to the poor fluidity and insufficient thermal stability of grease, the COF of base grease fluctuates violently and it increases obviously when the temperature is higher than 100 • C. It can be dramatically decreased by adding the modified talc nanosheets at a low concentration of 0.25 wt.%. The optimized concentration is about 0.5 wt.% (Figure 6d), where the COF can be decreased by about 40%. For the unmodified talc-based grease, the lubrication performance is not good and the COF is as high as 0.18. It confirms that the modified talc has a better lubrication performance. The load conditions also have a significant influence on the lubrication performance (Figure 6e,f). It can be seen that under the load varying from 150 N to 450 N, the modified talc-based grease maintains a relatively stable average COF (0.11-0.13), whereas the average COF of base grease is always higher than 0.15. Although the lubrication performance fails when the load increases to 550 N, the modified talc-based grease keeps a relatively longer stability than that of base grease. tion of 0.25 wt.%. The optimized concentration is about 0.5 wt.% (Figure 6d), where the COF can be decreased by about 40%. For the unmodified talc-based grease, the lubrication performance is not good and the COF is as high as 0.18. It confirms that the modified talc has a better lubrication performance. The load conditions also have a significant influence on the lubrication performance (Figure 6e,f). It can be seen that under the load varying from 150 N to 450 N, the modified talc-based grease maintains a relatively stable average COF (0.11-0.13), whereas the average COF of base grease is always higher than 0.15. Although the lubrication performance fails when the load increases to 550 N, the modified talc-based grease keeps a relatively longer stability than that of base grease.
Anti-Wear Performance
After the friction tests, the wear scars on the steel disc were successively investigated to evaluate the anti-wear effect of modified talc nanosheets. It can be seen that the wear width and wear depth of the disc lubricated by modified talc-based grease are much smaller than those of base grease even at the high temperature of 150 • C (Figure 7a,b). The wear behavior for base greases is worse when the temperature is 175 • C.
talc-based protective film formed on the rubbing surface, the modified talc-based grease exhibits a significant lubrication performance [13,42]. Although the wear width and wear depth increase with the increase in load, the modified talc-based grease exhibits a better anti-wear performance under different loads. The wear width and wear depth can be reduced by 26% and 66% under the load of 450 N (maximum contact stress: 3.5 GPa) (Figure 7c,d). The morphologies of the wear scars on the disc are characterized with an optical microscopy ( Figure 8). The addition of modified talc can effectively reduce the geometric size of wear scar and improve the anti-wear performance of the grease. At a high temperature (125 °C ), several furrows appear on the surface of discs lubricated by base grease. For comparison, there is no visible furrow mark when the added talc nanosheets concentration is 0.5 wt.%. Thus, it is confirmed that the modified talc nanosheets are able to enhance the grease performance at high temperature. In addition, in the case of base grease under high load (3.5 GPa), deep and wide wear scars are observed on the steel disc and the depths of the wear scar for the lubrication of base grease reach 13.36 µ m (Figure 9ac). These results suggest that base grease could not provide good anti-wear performance, It is because temperature has a crucial effect on the viscosity of grease. At high temperature, the bearing capacity of grease film decreases, and results in severe direct contact of interfacial asperities, leading to increase in friction. However, because of the modified talc-based protective film formed on the rubbing surface, the modified talc-based grease exhibits a significant lubrication performance [13,42]. Although the wear width and wear depth increase with the increase in load, the modified talc-based grease exhibits a better anti-wear performance under different loads. The wear width and wear depth can be reduced by 26% and 66% under the load of 450 N (maximum contact stress: 3.5 GPa) (Figure 7c,d).
The morphologies of the wear scars on the disc are characterized with an optical microscopy ( Figure 8). The addition of modified talc can effectively reduce the geometric size of wear scar and improve the anti-wear performance of the grease. At a high temperature (125 • C), several furrows appear on the surface of discs lubricated by base grease. For comparison, there is no visible furrow mark when the added talc nanosheets concentration is 0.5 wt.%. Thus, it is confirmed that the modified talc nanosheets are able to enhance the grease performance at high temperature. In addition, in the case of base grease under high load (3.5 GPa), deep and wide wear scars are observed on the steel disc and the depths of the wear scar for the lubrication of base grease reach 13.36 µm (Figure 9a-c). These results suggest that base grease could not provide good anti-wear performance, and could not protect the friction pair well in practical applications. In comparison, the rubbing surface lubricated by the nanosheets is very smooth under the load of 150 N (2.5 GPa). Although the wear depth and wear width increase at a higher load (3.5 GPa), the anti-wear performance of the nanosheets-filled grease under a high load is much better than that of unfilled grease. These results indicate that the modified talc nanosheet could provide an excellent anti-wear performance (Figure 9d-f). and could not protect the friction pair well in practical applications. In comparison, the rubbing surface lubricated by the nanosheets is very smooth under the load of 150 N (2.5 GPa). Although the wear depth and wear width increase at a higher load (3.5 GPa), the anti-wear performance of the nanosheets-filled grease under a high load is much better than that of unfilled grease. These results indicate that the modified talc nanosheet could provide an excellent anti-wear performance (Figure 9d-f). and could not protect the friction pair well in practical applications. In comparison, the rubbing surface lubricated by the nanosheets is very smooth under the load of 150 N (2.5 GPa). Although the wear depth and wear width increase at a higher load (3.5 GPa), the anti-wear performance of the nanosheets-filled grease under a high load is much better than that of unfilled grease. These results indicate that the modified talc nanosheet could provide an excellent anti-wear performance (Figure 9d-f). The typical SEM morphologies of the worn surfaces lubricated by base grease exhibits long wear tracks, as shown in Figure 10a, because worn surface asperities between the friction pairs could directly contact with each other, and scratch the friction interfaces during the boundary lubrication. The worn surface shows obvious abrasive and delamination wear (Figure 10c). Figure 10b,d show some slight scratches and tracks on the surface lubricated by the nanosheet-based grease. The above results agree with the result of optical morphologies of wear tracks in Figure 9. It can be seen that there are main Fe, Cr, C and O elements in or out of the wear area under the lubrication of base grease. For comparison, there is some residual talc on the cleaned surface, which is identified according to obvious Mg and Si elements appearing on the talc-lubricating surface, i.e., it results from the talc (Mg 3 Si 4 O 10 (OH) 2 ) adsorbed on the rubbing surface. It means that the talc is able to form a protective boundary film and repair the rubbing surface to enhance the lubrication performance. Easy layer-shearing effects and the thermostability of talc further improve the lubrication performance [31,32]. On the basis of the tribological analysis described above, as well as the convenient and economic modification route of the talc, the modified talc is confirmed as a promising additive for efficient and green lubrication. boundary lubrication. The worn surface shows obvious abrasive and delamination wear ( ure 10c). Figure 10b, d show some slight scratches and tracks on the surface lubricated by nanosheet-based grease. The above results agree with the result of optical morphologie wear tracks in Figure 9. It can be seen that there are main Fe, Cr, C and O elements in or of the wear area under the lubrication of base grease. For comparison, there is some resid talc on the cleaned surface, which is identified according to obvious Mg and Si elements pearing on the talc-lubricating surface, i.e., it results from the talc (Mg3Si4O10(OH)2) adsor on the rubbing surface. It means that the talc is able to form a protective boundary film repair the rubbing surface to enhance the lubrication performance. Easy layer-shearing eff and the thermostability of talc further improve the lubrication performance [31,32]. On basis of the tribological analysis described above, as well as the convenient and econo modification route of the talc, the modified talc is confirmed as a promising additive for cient and green lubrication.
Conclusions
The 2D talc nanosheets were modified by a silane coupling agent with the assistance of highly-energy ball milling. The effect of talc nanosheets on the lubrication performance of grease was studied using a SRV-4 tribometer. The micromorphology of the modified talc was characterized by SEM, EDS, and TEM methods, together with XRD, FTIR and TG. In addition, the lubrication performance of grease was investigated in detail. The conclusions achieved are as follows: (1) The silane coupling agent was successfully coupled with the talc nanosheets via ball milling. The modified talc nanosheet has a non-defected crystal structure with a thickness and size of about 20 nm and 2 µm. (2) The modified talc has a much better lubrication performance than the non-modified talc and base grease. The optimum addition level of modified talc is 0.5 wt.%. The modified talc can greatly enhance the lubrication performance under high temperature and high load. Compared with the base grease, the coefficient of friction and wear depth can be reduced by 40% and 66% at high temperature (150 • C) and high load (3.5 GPa), respectively. (3) The modified talc nanosheets via ball milling together with grafting of the agent can be uniformly dispersed in viscous grease, which has a good chance to make the nanosheets enter the wear interface quickly and form a protective adsorbed tribofilm. This study provides a green and economical way for using nanomaterial as an efficient lubricant additive and has great potential for application. | 6,235.6 | 2022-05-01T00:00:00.000 | [
"Materials Science",
"Engineering"
] |
Interaction of POB1, a downstream molecule of small G protein Ral, with PAG2, a paxillin-binding protein, is involved in cell migration.
POB1 was previously identified as a RalBP1-binding protein. POB1 and RalBP1 function downstream of small G protein Ral and regulate receptor-mediated endocytosis. To look for additional functions of POB1, we screened for POB1-binding proteins using a yeast two-hybrid method and found that POB1 interacts with mouse ASAP1, which is a human PAG2 homolog. PAG2 is a paxillin-associated protein with ADP-ribosylation factor GTPase-activating protein activity. POB1 formed a complex with PAG2 in intact cells. The carboxyl-terminal region containing the proline-rich motifs of POB1 directly bound to the carboxyl-terminal region including the SH3 domain of PAG2. Substitutions of Pro(423) and Pro(426) with Ala (POB1(PA)) impaired the binding of POB1 to PAG2. Expression of PAG2 inhibited fibronectin-dependent migration and paxillin recruitment to focal contacts of CHO-IR cells. Co-expression with POB1 but not with POB1(PA) suppressed the inhibitory action of PAG2 on cell migration and paxillin localization. These results suggest that POB1 interacts with PAG2 through its proline-rich motif, thereby regulating cell migration.
POB1 (5). Eps15 binds to ␣-adaptin, a subunit of the clathrin adaptor complex, AP-2 (6). The AP-2-binding site of Eps15 acts dominant negatively in blocking the endocytosis of EGF, transferrin, or Sindbis virus, indicating that Eps15 is actively required for the endocytic process (7,8). Eps15R has a 47% amino acid identity with and exhibits similar characteristics to Eps15 (4,9). The POB1-related protein Reps1 has been identified as a RalBP1-binding protein (10). Intersectin (Ese) has five SH3 domains in addition to two EH domains and is involved in the regulation of internalization of the transferrin receptor (11,12). Pan1 and End3 are Saccharomyces cerevisiae dimeric partners that are necessary for endocytosis of the ␣-mating factor receptor and for normal organization of the actin cytoskeleton (13,14). Thus, the EH domain-containing proteins regulate endocytosis.
EGF and insulin stimulate the GDP/GTP exchange of Ral through Ras and RalGDS (15)(16)(17), and the GTP-bound active form of Ral interacts with RalBP1 (1). The carboxyl-terminal region of POB1 binds to RalBP1 (2). Because the binding sites of Ral and POB1 on RalBP1 are different, these three proteins form a ternary complex. EGF stimulates tyrosine phosphorylation of POB1 and induces the complex formation between EGF receptor and POB1 (2). The EH domain of POB1 associates with Eps15 and Epsin (18,19). Epsin also regulates endocytosis by directly binding to phospholipids (20), ␣-adaptin (21), and clathrin (22,23). Expression of the EH domain or the carboxyl-terminal region of POB1 inhibits the internalization of EGF and insulin (18). Therefore, it is conceivable that Ral, RalBP1, and POB1 regulate receptor-mediated endocytosis by transmitting the signal from receptors to Eps15 and Epsin.
The Arf family of small G proteins is divided into three classes based largely on sequence similarity: class I (Arfs 1-3), class II (Arfs 4 and 5), and class III (Arf6) (24). By linking GTP binding and hydrolysis, Arfs regulate membrane trafficking at various steps (25,26). For instance, Arf6 has been implicated in the regulation of membrane trafficking between the plasma membrane and a specialized endocytic component. Moreover, its function has been linked to cytoskeletal reorganization (27,28). As with other small G proteins, the activity of Arfs is regulated by guanine nucleotide exchange factors and GAPs. It has been shown that ArfGAP is involved in regulating the organization of focal adhesions (29,30). Evidence for a direct link between Arf signaling and focal adhesions came initially from the identification of the ArfGAP protein as a paxillinbinding protein. There are several ArfGAP families. All Arf-GAP proteins share homology within the zinc-finger-containing ArfGAP domain and ankyrin repeat region. Among ArfGAP family proteins, PAG3 contains a pleckstrin homology domain in an extended amino terminus and has a proline-rich sequence followed by an SH3 domain at the carboxyl terminus (31). PAG3 binds to paxillin and serves as a GAP for Arf6. Overexpression of PAG3 in fibroblasts inhibits cell motility and reduces the paxillin recruitment to focal contacts in a GAP-dependent manner (31). These results suggest that PAG3 plays a role in mediating changes in cell motility.
To find additional functions of POB1, we screened proteins that bind to POB1. Here we report that the proline-rich domain of POB1 interacts with the SH3 domain of PAG2, a PAG3 homolog. Furthermore, we show that the functional interaction of POB1 with PAG2 may regulate cell migration. These results suggest that POB1 and PAG2 link the processes of endocytosis and cell motility.
EXPERIMENTAL PROCEDURES
Materials and Chemicals-Recombinant baculovirus expressing GST-POB1 was provided by Dr. Y. Matsuura (Research Center for Emerging Infectious Diseases, Research Institute for Microbial Diseases, Osaka University). Hygromycin-resistant CHO-IR cells that stably express POB1, PAG2, or their mutants were propagated as described (32). CHO-IR cells stably expressing both POB1 and PAG2, both POB1 P423A/P426A and PAG2, or both POB1 and PAG2-(1-703) were generated by selecting with Blasticidin. The rabbit polyclonal anti-GST and anti-MBP antibodies were made by a standard method. The rabbit polyclonal anti-POB1 and anti-PAG2 antibodies were prepared by immunization with recombinant POB1-(322-521) and GST-PAG2-(935-1002), respectively. The mouse monoclonal anti-HA antibody 12CA5 was kindly provided by Dr. Q. Hu (Chiron Corp., Emeryville, CA). The rabbit polyclonal anti-PAG3 antibody was prepared as described (31). The mouse monoclonal anti-Myc antibody was prepared from 9E10 cells. GST and MBP fusion proteins were purified from Escherichia coli according to the manufacturer's instructions. GST-POB1 was purified from Spodoptera frugiperda (Sf) 9 cells. Other materials were from commercial sources.
Two-hybrid Screening-Yeast strain Y190 was used as a host for the two-hybrid screening (CLONTECH Laboratories Inc., Palo Alto, CA). Yeast cells were grown on rich medium (YAPD) containing 2% glucose, 2% Bact-peptone, 1% Bact-yeast extract, and 0.002% adenine sulfate. Yeast transformations were performed by the lithium acetate method. Transformants were selected on SD medium containing 2% glucose, 0.67% yeast nitrogen base without amino acids, and necessary supplements. Y190 strain carrying pGBKT7/POB1-(321-521), in which POB1-(321-521) was expressed as a fusion protein with the GAL4 DNAbinding domain, was transformed with a mouse brain cDNA library constructed in pACT2, in which cDNA was expressed as a fusion protein with the GAL4 activator domain. Approximately 1.6 ϫ 10 6 transformants were screened for the growth on SD plate medium lacking tryptophan, leucine, and histidine as evidenced by transactivation of a GAL4-HIS3 reporter gene and histidine prototrophy. His ϩ colonies were scored for -galactosidase activity. Plasmids harboring cDNAs were recovered from positive colonies and introduced by electroporation into E. coli HB101 on M9 plates lacking leucine. HB101 is leuB Ϫ , and this defect can be complemented by the LEU2 gene in the library plasmids. The library plasmids were then recovered from HB101 and transformed into Y190 containing pGBKT7/POB1-(321-521). The nucleotide sequence of the plasmid cDNAs, which conferred the LacZ ϩ phenotype on Y190 containing pGBKT7/POB1-(321-521), was determined.
To examine whether paxillin interacts with the complex of PAG2 and POB1, CHO-IR cells (10-cm diameter plate) expressing HA-POB1 and/or GFP-PAG2 were lysed in 0.25 ml of lysis buffer. The lysates (480 g of protein) were incubated with 5 g of GST-paxillin ␣ immobilized on glutathione-Sepharose 4B for 1 h at 4°C. After glutathione-Sepharose 4B was precipitated by centrifugation, the precipitates were probed with the anti-GFP and anti-HA antibodies.
When the complex formation of POB1 with PAG3 was examined, 293 cells (6-cm diameter dishes) were lysed in 200 l of lysis buffer. The lysates (180 g of protein) were immunoprecipitated with the anti-POB1 antibody, and the immunoprecipitates were probed with the anti-PAG3 and anti-POB1 antibodies.
Direct Binding of POB1 and Paxillin to PAG2 in Vitro-Various GST-fused POB1 deletion mutants (0.5 M each) were incubated with 20 pmol of MBP-PAG2-(1002-1132) immobilized on amylose resin in 100 l of reaction mixture (20 mM Tris/HCl, pH 7.5, and 1 mM dithiothreitol) for 1 h at 4°C. After the resin was precipitated by centrifugation, the precipitates were probed with the anti-GST antibody.
To show the simultaneous binding of PAG2 to POB1 and paxillin in vitro, 0.5 M GST-POB1-(322-521) and/or 1-4 M GST-paxillin were incubated with 20 pmol of MBP-PAG2-(1002-1132) or MBP immobilized on amylose resin in 100 l of reaction mixture (20 mM Tris/HCl, pH 7.5, and 1 mM dithiothreitol) for 1 h at 4°C. After the resin was precipitated by centrifugation, the precipitates were probed with the anti-GST antibody.
SPR Spectroscopy-The binding of MBP-PAG2-(1002-1132) to GST-POB1 was investigated by real-time SPR spectroscopy (BIACORE X, BIAcore System, Uppsala, Sweden). Measurements were performed at 25°C in HBS-EP buffer (10 mM HEPES/NaOH, pH 7.4, 150 mM NaCl, 3 mM EDTA, 0.005% polysorbate 20) at a flow rate of 20 l/min. GST-POB1 was coupled to a CM5 sensor tip in amounts to yield 800 resonance units (RU) (1 RU is equal to 1 pg/mm 2 ), via the goat anti-GST antibody that was covalently coupled to the surface of the sensor tip by standard amine-coupling chemistry. GST was coupled to the reference cell. MBP-PAG2-(1002-1132) was injected for 180 s followed by elution in HBS-EP buffer for 180 s. The observed changes in the relative diffraction indices, which represents the mass on the sensor tip surface, were recorded as a function of time. The value obtained from the reference cell, which represented nonspecific binding of MBP fusion proteins to GST, was subtracted from that with GST-POB1. Association and dissociation constants of the PAG2-POB1 complex were calculated using the BIAevaluation program version 3.1 (BIAcore). K d was calculated by fitting the data to the equation Cell Adhesion Assay-The CHO-IR cells (5 ϫ 10 4 ) were added to a 96-well dish that was precoated with 10 g/ml fibronectin or BSA in PBS for 2 h. After 30 min of incubation, the cells were washed with PBS three times. Adherent cells were then fixed and stained with 0.1% crystal violet in 20% methanol for 5 min at room temperature and were washed with PBS extensively. The stain was eluted with 100 l of 50% ethanol, and the absorbance at 590 nm was measured as described (37).
Cell Migration Assay-The cell migration assay was performed using a modified Boyden chamber (tissue culture treated, 6.5-mm diameter, 10-m thickness, 8-m pore; Transwell, Costar, Cambridge, MA) as described (38). In brief, only the underside surface of the polycarbonate membrane on the upper chamber was coated with 10 g/ml fibronectin or BSA in PBS for 2 h. After the chamber was rinsed with PBS, it was placed into the lower chamber filled with 400 l of Ham's F-12 medium containing 1% fetal calf serum. CHO-IR cells (2.5 ϫ 10 4 ) suspended in the Ham's F-12 medium containing 0.1% BSA at 2.5 ϫ 10 5 cells/ml were applied to the upper chamber and allowed to migrate to the underside of the upper chamber for 3 h at 37°C with 5% CO 2 . After the nonmigrated cells on the upper membrane surface were removed with a cotton swab, cells that migrated to the underside of the upper chamber were fixed with 4% paraformaldehyde in PBS and were stained with propidium-iodide solution. The number of the stained cells was counted and percent cell migration was calculated by dividing the number of stained cells by the number of applied cells.
Immunofluorescence Study-The CHO-IR cells expressing GST-PAG2, GFP-POB1, or GFP-POB1(PA) were grown on coverslips and then fixed for 20 min in PBS containing 4% paraformaldehyde. The cells were washed with PBS three times and then permeabilized with PBS containing 0.1% Triton X-100 and 2 mg/ml BSA for 20 min. They were washed and incubated for 1 h with the mouse monoclonal anti-paxillin and rabbit polyclonal anti-GST antibodies. After being washed with PBS, the cells were further incubated for 1 h with Cy5-labeled antimouse IgG and Cy3-labeled anti-rabbit IgG. The coverslips were washed with PBS and mounted on glass slides, and the fluorescence of Cy5, Cy3, and GFP was viewed with a confocal laser-scanning microscope (LSM510, Carl-Zeiss, Jena, Germany). To determine the effects of POB1 on the inhibitory action of PAG2 for the paxillin recruitment to focal contacts, the number of cells where paxillin was recruited to focal contacts was divided by the total number of the cells counted. Approximately 100 -300 cells were examined in each experiment.
Identification of PAG2 as a POB1-binding Protein-Various
constructs used in this study are shown in Fig. 1. To discover new functions of POB1, we attempted to identify POB1-binding protein(s). We screened a mouse brain cDNA library by the yeast two-hybrid method using the carboxyl-terminal region of POB1 (POB1-(322-521)) as bait. Several clones were found to confer both His ϩ and LacZ ϩ phenotypes, and three of them overlapped, encoding the carboxyl-terminal region of mouse ASAP1 (ASAP1-(1050 -1147)). ASAP1 has been identified as an ArfGAP and contains a zinc finger domain similar to that required for GAP activity for Arf (39). ASAP1 also contains a number of domains that are likely to be involved in regulation and/or localization: a pleckstrin homology (PH) domain, three ankyrin (ANK) repeats, a proline-rich region, and an SH3 domain. To examine whether POB1 forms a complex with ASAP1 in intact cells, we expressed HA-POB1 and Myc-ASAP1-(1050 -1147) in COS cells. When the lysates were immunoprecipitated with the anti-Myc antibody, HA-POB1 was detected in the Myc-ASAP1-(1050 -1147) immune complexes under the conditions that HA-POB1 formed a complex with Myc-RalBP1 ( Fig. 2A). Previously we isolated PAG2 and PAG3 as paxillin-binding proteins (31). Because human PAG2 has a higher homology with mouse ASAP1 (95% identity) than human PAG3 does (58% identity), we examined whether POB1 interacts with full-length PAG2 in intact cells. GST-PAG2 was expressed in CHO-IR cells stably expressing HA-POB1 (Fig. 2B, lanes 1 and 2). When the lysates of CHO-IR cells expressing both GST-PAG2 and HA-POB1 were precipitated with glutathione-Sepharose, HA-POB1 was co-precipitated with GST-PAG2 (Fig. 2B, lanes 3 and 4).
Next, we asked whether endogenous PAG2 associated with endogenous POB1 in CHO-IR cells. When the lysates of CHO-IR cells were immunoprecipitated with the anti-POB1 antibody, PAG2 was detected in the POB1 immune complex (Fig. 2C, lanes 1-3). We also examined whether endogenous PAG3 interacted with endogenous POB1 in intact cells. Because PAG3 was expressed only slightly in CHO-IR cells, we used 293 cells. Endogenous PAG3 was observed in the POB1 immune complex from 293 cells (Fig. 2C, lanes 4 -6). These results indicate that POB1 forms a complex with PAG2 and PAG3 in intact cells at endogenous levels.
To determine the association and dissociation rates of complex formation between POB1 and PAG2, we performed real time SPR analysis (Fig. 3B). For this purpose, GST-POB1 and GST were immobilized onto a CM5 sensor tip surface via the anti-GST antibody, which was covalently coupled to the surface by standard amine-coupling chemistry. GST was immobilized onto the surface of the reference cell on the same sensor tip. Four different concentrations (17,33,67, and 133 nM) of MBP-PAG2-(1002-1132) or 100 nM MBP were injected onto the surface of the sensor tip at 25°C for 180 s to form the complex, and then the sensor tip was washed with the buffer for 180 s to dissociate the complex. The association rate k a for the binding of MBP-PAG2-(1002-1132) was determined to be 1.06 Ϯ 0.02 ϫ 10 5 M Ϫ1 s Ϫ1 and the dissociation rate k d to be 1.44 Ϯ 0.02 ϫ 10 Ϫ3 s Ϫ1 . The static dissociation constant K d was calculated as 13.6 nM. Thus, POB1 binds to PAG2 with high affinity, consistent with the observations that these proteins form a complex at endogenous levels.
Under the same conditions, both GST-POB1-(322-521) and GST-POB1-(322-521)(PA) interacted with MBP-RalBP1-(364 -647), which is known to contain the POB1-binding region (Fig. 4B, lanes 5 and 6). These observations were also confirmed in intact cells. Although Myc-RalBP1 formed a complex with HA-POB1(PA) in CHO-IR cells, GST-PAG2 did not (Fig. 4C). These results clearly indicate that Pro 423 and Pro 426 of POB1 are essential for the interaction with PAG2 and that the sites on POB1 that bind to PAG2 and RalBP1 are different.
Functional Interaction of POB1 with PAG2
immunoprecipitated with the anti-Myc antibody, GST-PAG2 was faintly detected in the Myc-RalBP1 immune complex, suggesting that RalBP1 forms a complex with PAG2 via endogenous POB1 (Fig. 5B, lanes 2 and 6). Additional expression of HA-POB1 enhanced the complex formation of GST-PAG2 and Myc-RalBP1 (Fig. 5B, lanes 3 and 7). However, HA-POB1(PA) did not allow Myc-RalBP1 to associate with GST-PAG2, and instead inhibited formation of the complex (Fig. 5B, lanes 4 and 8). These results suggest that RalBP1, POB1, and PAG2 may form a ternary complex.
Complex Formation of POB1, PAG2, and Paxillin-Because PAG2 is a homolog of PAG3, we examined whether PAG2 also binds to paxillin. When the lysates of CHO-IR cells expressing either GFP-PAG2 or HA-POB1 were incubated with GST-paxillin, GFP-PAG2 associated with GST-paxillin, but HA-POB1 interacted with GST-paxillin only faintly (Fig. 6A, lanes 7-10). When the lysates of CHO-IR cells expressing both GFP-PAG2 and HA-POB1 were incubated with GST-paxillin, both proteins formed a complex with GST-paxillin but not with GST (Fig. 6A, lanes 11 and 12). HA-POB1 formed a complex with GST-paxillin more efficiently than when HA-POB1 was expressed alone (Fig. 6A, lanes 10 and 12). These results suggest that POB1 forms a complex with paxillin through PAG2 and that the complex of POB1 and PAG2 can bind to paxillin.
Effects of POB1 on Cell Adhesion and Migration-Cell adhesion and migratory activities are primarily mediated by inte- 1 and 2). The same lysates were precipitated with glutathione-Sepharose, and the precipitates were probed with the anti-GST and anti-HA antibodies (lanes 3 and 4). C, complex formation of POB1 with PAG2 or PAG3 at endogenous level. The lysates of CHO-IR cells (lanes grin adhesion to the extracellular matrix. As it was shown that overexpression of PAG3, a homolog of PAG2, decreases cell migratory activity (31), we examined whether POB1 and PAG2 affect these activities of CHO-IR cells. To this end, we generated CHO-IR cells stably expressing POB1 mutants and/or PAG2 mutants (Fig. 7A). The cell adhesiveness toward fibronectin of CHO-IR cells stably expressing HA-POB1 (CHO-IR/POB1) was similar to that of CHO-IR cells (Fig. 7B). Furthermore, expression of HA-POB1(PA), GST-PAG2, GST-PAG2N, GST-PAG2C, HA-POB1 and GST-PAG2, HA-POB1(PA) and GST-PAG2, or HA-POB1 and GST-PAG2N did not affect cell adhesiveness (Fig. 7B). Therefore, it is likely that POB1 and PAG2 are not involved in cell adhesion. Overexpression of POB1 or POB1(PA) in CHO-IR cells did not affect the cell migratory activity on fibronectin (Fig. 7C). PAG2 caused a several-fold decrease in the cell migratory activity (Fig. 7C). It has been shown that the ArfGAP activity of PAG3 is essential for its ability to suppress cell migration (31). The amino-terminal region of PAG2 contains the ArfGAP domain (PAG2N-(1-703)). CHO-IR/PAG2N-(1-703) decreased cell migration activity, whereas CHO-IR/PAG2C-(704 -1132) showed similar activity to CHO-IR cells. These results suggest that the ArfGAP domain, but not the POB1-binding domain, of PAG2 is important for the ability to suppress cell migration. Co-expression with POB1 but not with POB1(PA) suppressed the PAG2-induced inhibition of motility (Fig. 7C). Moreover, co-expression with POB1 could not suppress the inhibition of motility induced by PAG2N-(1-703) (Fig. 7C). These results suggest that the interaction of POB1 with PAG2 prevents PAG2 from inhibiting cell migration.
Effects of POB1 on the Inhibitory Action of PAG2 on the Paxillin Recruitment to Focal Contacts-Paxillin was condensed at focal adhesion plaques at the bottom of the cells (Fig. 8A). Overexpression of PAG3 caused loss of endogenous paxillin recruitment to focal contacts (31). As in the case of PAG3, when GST-PAG2 was overexpressed in CHO-IR cells, the staining of paxillin decreased (Fig. 8A, a). When GST-PAG2 was co-expressed with GFP-POB1, paxillin was observed as punctate structures, showing similar staining as in the surrounding normal cells (Fig. 8A, d). GFP-POB1(PA) did not influence the effect of GST-PAG2 on paxillin staining (Fig. 8A, g). These results suggest that POB1 suppresses the inhibitory action of PAG2 on the paxillin recruitment to focal contacts by binding to PAG2. To quantitatively determine the effects of POB1 on the inhibitory action of PAG2 on the paxillin recruitment, the number of cells that showed clear staining of paxillin at the cell bottom in the same slice was counted (Fig. 8B). In the normal cells surrounding the cells expressing GST-PAG2 and/or GFP-POB1, paxillin was observed in 68 Ϯ 5.6% of all the cells examined. The percentages of the cells with clear staining of paxillin at the cell bottom in the cells expressing GST-PAG2 alone, GST-PAG2 and GFP-POB1, and GST-PAG2 and GFP-POB1(PA) were 24 Ϯ 1.6%, 68 Ϯ 11.1%, and 39 Ϯ 9.8%, respectively.
DISCUSSION
In this study, we demonstrated the interaction of POB1 with PAG2. Endogenous PAG2 was detected in the endogenous POB1 immune complex from CHO-IR cells. Sf9-cell-produced GST-POB1 bound to bacterial cell-produced MBP-PAG2-(1002-1132) containing the SH3 domain with a K d value of 13.6 nM. Therefore, it is conceivable that POB1 binds directly to PAG2 under physiological conditions. Furthermore, we demonstrated that the proline-rich motif of POB1 is essential for the binding of POB1 to PAG2. POB1 has three proline-rich motifs, PPTPPPRP 345 , PPPPALPPRP 383 , and PPSKPIR 428 . It is generally thought that the proline-rich motifs bind to several proteins such as profilin and to the EVH1, SH3, and WW domains (40). The core motif that binds to the SH3 domain is PXXP, and this motif is further classified into class I and class II. The class I motif is (R/K)XXPXXP, which binds to the SH3 domains of Src, Abl, Fyn, and Lyn. The class II motif is PXXPX(R/K), which binds to the SH3 domains of Grb2, Nck, and Crk. All of the proline-rich motifs of POB1 are class II of the SH3-domain binding motifs. Because substitution of two proline residues with alanine in the third motif of POB1 impaired its binding to PAG2, the third proline motif is essential for the binding to PAG2. These results suggest that the proline-rich motifs of POB1 interact with the SH3 domain of PAG2. It has been shown that the SH3 domain of PAG3/PAP␣ binds to Pyk2 and that activation of Pyk2 leads to tyrosine phosphorylation of PAP␣ (31,41). Because the SH3 domain of PAG2 shares 76% identity with that of PAG3, PAG2 may interact with Pyk2. It remains to be clarified whether POB1 affects the interaction of PAG2 with Pyk2. Previously we showed that among the SH3domain-containing proteins, Grb2 but not Nck and Crk binds to POB1 (2). Because Grb2 bound to both POB1 and POB1(PA) (data not shown), it seems that the third proline-rich motif of POB1 is not essential for its binding to Grb2, suggesting that PAG2 and Grb2 bind to different sites of POB1.
Several lines of evidence indicate that ArfGAP family members, including GIT, PAG3/PAP␣, and ASAP1/DEF-1, regulate actin cytoskeletal dynamics (30). PAG3 interacts with paxillin, which acts as an adaptor molecule in integrin signaling and is localized to focal contacts (29,31). PAG3 is diffusely distributed in the cytoplasm in premature monocytes but becomes localized at cell periphery in mature monocytes (31). However, PAG3 does not accumulate at focal contacts, suggesting that PAG3 is not an integrin assembly protein. Overexpression of PAG3 in COS-7 and U937 cells causes a loss of the paxillin recruitment to focal adhesions and inhibits cell motility in a GAP-dependent manner. Overexpression of PAG2 also impaired cell migratory activities and inhibited paxillin recruitment to focal contacts. This does not always reflect that PAG2 negatively regulates cell migration because overexpression of PAG2 may interfere with the functions of proteins that are involved in the cell migration through recruitment of the binding partners even though PAG2 is a positive regulator.
The amino-terminal region of PAG2 containing the ArfGAP domain, but not the carboxyl-terminal region containing the binding sites of paxillin and POB1, inhibited cell migration. Taken together with the observations that the activities of Arfs are involved in the focal adhesion recruitment of paxillin (31,42), it is conceivable that the ArfGAP activity is essential for these activities of PAG2, but we do not know the physiological roles of the paxillin-binding activity of PAG2 for them. We also showed that POB1, but not POB1(PA), restores cell motility and the paxillin recruitment to focal contacts, which are inhibited by PAG2. Moreover, POB1 could not rescue the inhibition by the amino-terminal region of PAG2 that lacks the POB1binding site. Therefore, the interaction of POB1 with PAG2 may regulate the paxillin recruitment to focal contacts, but we do not know the mechanism at present. One possibility might be that POB1 participates in the recruitment of PAG2 to proper subcellular areas in which PAG2 may act as a GAP for Arfs, resulting in the regulation of the subcellular positioning of paxillin. It has been proposed that primer proteins including the coatmer and the GTP-bound form of Arf at the membranes of the endoplasmic reticulum and the Golgi apparatus influence the catalytic activity of ArfGAP1 (43,44). Therefore, complex formation between paxillin, PAG2, and POB1 at certain subcellular areas may constitute a signal necessary for the onset or the enhancement of the catalytic GAP activity of PAG2 toward the GTP-bound form of Arfs.
Cell locomotion is driven by protrusive activity at the leading edge of the cell where continuous remodeling of actin cytoskelton and adhesive contacts is required (45). Endocytosed membrane is reinserted at the leading edge of migrating cells, extending the front of the cell forward. For instance, recycling transferrin receptors and low density lipoprotein receptors are distributed to the cell front of migrating fibroblasts and Racinduced ruffles (46,47). Therefore, it is likely that the random reinsertion of internalized membranes at the surface of a resting cell is redirected to the site of protrusion when migration is induced by mitogenic stimuli. Arf6 is implicated in the regulation of membrane trafficking between the recycling endosomal compartment and the plasma membrane, based on the specific localization of Arf6 in these compartments and the effects of its overexpression on transferrin uptake and recycling to the cell surface (48,49). Arf6 co-localizes with Rac1, which is involved in the formation of actin-rich ruffles and lamellipodia, at the plasma membrane and on recycling endosomes (50). Moreover, the ArfGAP family proteins interact with proteins involved in both cell adhesion and actin organization (30). Therefore, it has been speculated that ArfGAP is involved in the regulation of Arf-mediated membrane recycling and protrusion during cell locomotion.
We have demonstrated that small G protein Ral and its downstream molecules, RalBP1 and POB1, are involved in receptor-mediated endocytosis of EGF and insulin (18). Furthermore, we have found that Eps15 and Epsin bind directly to the EH domain of POB1 (18,19). These results suggest that the signaling from Ral to Eps15 and Epsin through RalBP1 and POB1 regulates receptor-mediated endocytosis. Because Eps15 and Epsin are core proteins that regulate endocytosis, POB1 may be able to link endocytosis and cell migration. The binding sites of POB1 for Epsin, RalBP1, and PAG2 are different. Taken together with the observations that Ral regulates both actin cytoskeletal remodeling and vesicle transport (18,37,51,52), it is intriguing to speculate that POB1 may function as a scaffold protein in that it interacts with proteins involved in endocytosis and migration to create a multi-protein complex. Further analysis would be necessary to understand how these complex interactions are temporally and spacially coordinated during cell migration. | 6,328 | 2002-10-11T00:00:00.000 | [
"Biology"
] |
Parent Hamiltonian Reconstruction of Jastrow-Gutzwiller Wavefunctions
Variational wave functions have been a successful tool to investigate the properties of quantum spin liquids. Finding their parent Hamiltonians is of primary interest for the experimental simulation of these strongly correlated phases, and for gathering additional insights on their stability. In this work, we systematically reconstruct approximate spin-chain parent Hamiltonians for Jastrow-Gutzwiller wave functions, which share several features with quantum spin liquid wave-functions in two dimensions. Firstly, we determine the different phases encoded in the parameter space through their correlation functions and entanglement content. Secondly, we apply a recently proposed entanglement-guided method to reconstruct parent Hamiltonians to these states, which constrains the search to operators describing relativistic low-energy field theories - as expected for deconfined phases of gauge theories relevant to quantum spin liquids. The quality of the results is discussed using different quantities and comparing to exactly known parent Hamiltonians at specific points in parameter space. Our findings provide guiding principles for experimental Hamiltonian engineering of this class of states.
Introduction
Variational wave functions play a key role in the understanding of quantum phases of matter [1][2][3][4][5][6][7][8]. A paradigmatic example is Laughlin wave functions [5], which can be formulated as parametric Jastrow states reproducing several key features of certain fractional quantum Hall effects [9]. Shortly after this, resonating valence bond (RVB) states have been employed as effective descriptions of high-temperature superconductors [6,7,10], and later on, have been linked to fractional quantum Hall physics in Ref. [8]. These early successes boosted variational wave functions as theoretical tools to provide simple pictures for a variety of quantum phases, including topological matter, low-dimensional systems, and tensor networks [11][12][13][14][15].
Perhaps, among these applications, one of the most fruitful has been in the field of quantum spin liquids [16][17][18][19][20][21][22]. These are quantum phases characterized by strong correlations and longrange entanglement among arbitrary far subregions of the system [23], and for these reasons, semi-classical pictures fail in describing the phenomena involved. Variational wave functions have been used to distill generic properties such as correlation functions and entanglement [14].
Interestingly, despite the conceptual simplicity of Jastrow wave functions, it is often challenging to find the corresponding parent Hamiltonians -that is, the Hamiltonians supporting these wave functions as ground states. The major obstruction is that, given a Hamiltonian on a lattice (possibly with frustration terms), quantum fluctuations may cooperate and induce an ordered ground state. This phenomenon is typically referred to as "order-by-disorder" [13]. This problem is of primary importance also due to the latest experimental breakthrough in quantum engineering of synthetic systems [24][25][26][27][28]. In fact, the high degree of interaction tunability of these platforms offers new perspectives and possibilities in otherwise hardly achievable phases of matter, including spin liquids, once parent Hamiltonians are (approximately) identified.
Most of the works in parent Hamiltonian construction studied specific variational states using insightful analytic manipulations [29][30][31][32][33][34][35][36][37][38][39][40][41][42][43][44]. Very recently, a series of novel techniques based on systematic approaches have been considered in Ref. [45][46][47][48][49][50]. Indeed, the authors of the latter works introduced new efficient computational algorithms, which remarkably scale polynomially in the system size when restricting the search to local Hamiltonians that have a given initial state as the input eigenstate. To benchmark their techniques, they considered the ground state of some a priori known Hamiltonian as input and checked if the output reconstructed operator coincided with that Hamiltonian. So far, however, there have been no applications of such methods to generic spin liquid variational wave functions, whose parent Hamiltonians are still undetermined.
The present work is the first step in this direction. For concreteness, here we study the class of 1D Jastrow-Gutzwiller variational wave functions [30,51]. These states share two key features with their two-dimensional cousins employed as effective descriptions of quantum spin liquids: they describe extensive superpositions over some (spatially local) state basis, and they have in general as weights analytic functions of the space coordinates. Despite their common appearance, their parent Hamiltonians are not known except for a few fine-tuned cases, amenable to exact solutions. We use an entanglement-guided algorithm presented in Ref. [50] to search local parent Hamiltonians for these states. This method relies on the Bisognano Wichmann theorem [52,53], a quantum field theory result that links systematically the local Hamiltonian density to its ground state reduced density matrix. Its advantage with respect to the other above-mentioned techniques resides in certifying the input state as the ground state of the reconstructed parent Hamiltonian. Indeed, although the methods in Ref. [45][46][47][48][49] are of broader applicability (for instance, they allow for extensions to timedependent problems), they typically certify the ansatz state to be a generic eigenstate, and not the ground state, of the output operator. The main disadvantage is that the method is not applicable in case the wave function cannot be cast as the ground state of Hamiltonian operator supporting low-energy relativistic excitations.
Since the Bisognano-Wichmann technique requires the input state to exhibits relativistic low lying physics, we first investigate the entanglement and correlation properties of these wave functions, identifying a region where the algorithm is expected to perform better. In this regime, we obtain local approximate parent Hamiltonian searching through different algebras of local operators. To check our results, we computed the relative entropy, the correlation functions and the overlap between their ground state and the Jastrow-Gutzwiller wave functions, obtaining fidelities ranging between 95% to over 99%. In addition, we computed the relative error between the ground state energy and the Jastrow-Gutzwiller variational energy of the reconstructed Hamiltonian. In all the considered cases, the relative error is less than 1%, even in the extrapolated thermodynamic limit. We perform systematic searches by increasing both system sizes and interaction range. These results suggest that the exact, yet unknown, parent Hamiltonians of these states exhibit long-range features.
In addition, the method allows us to perform direct parent Hamiltonian searches utilizing simple long-range interactions in the form of monotonous power-law potentials. We find that, while considerably improving the parent Hamiltonian search, such simple long-range interactions are not always sufficiently rich to capture the (unapparent) complexity of Jastrow-Gutzwiller wave functions. These results indicate that the search for exact -albeit long-ranged -parent Hamiltonians for 2D Jastrow-Gutzwiller might be particularly challenging, a fact which is compatible with the scarcity of exact results in this context (with some notable exceptions, see Ref. [29,40]).
The remaining of this paper is structured as follows. In Section 2 we introduce the Jastrow-Gutzwiller states and discuss their physical content through participation spectrum, entanglement entropy and correlation functions. In Section 3 we summarize the Bisognano-Wichmann Ansatz method which we employ in Section 4 to reconstruct various parent Hamiltonians for the above-considered states. The last section is for conclusions and outlooks.
Model wave functions
The Jastrow-Gutzwiller (JG) wave functions are paradigmatic states appearing in several contexts, from integrability to topology (e.g. Laughlin states), to quantum spin liquids. They are characterized by an extensive superposition of spatially local states, and the local weights of the wave functions are captured by polynomials. Throughout this paper, we investigate the one-dimensional case defined on a periodic chain Λ of length L. This setting permits the understanding of finite-volume effects in a systematic manner, as well as enables comparison to exact results.
Let us introduce the wave functions of interest, through the variables n i ∈ {0, 1} defined at each site i ∈ Λ. In the basis {|n 1 n 2 . . . n L }, these states read: Here the sum is over combinations P N {n} constrained by i n i = N . Pictorially, the {n i } variables are occupation numbers of hard-core bosons living on the lattice. The real parameter α and the filling fraction ν = N/L control the properties of the states. For specific combined values of ν and α, conformal field theory calculations have been used to derive exact results pertaining the parent Hamiltonians of these states [42-44, 54, 55]. Throughout this paper, we will consider exclusively the half-filling case ν = 1/2 and L even; the main motivation being that, in spin language, this regime captures both paramagnetic and antiferromagnetic phenomenology. Within this setting, exact results are available only for α ∈ {0, 1, 2}. In Ref. [56], it was proven that α = 0 corresponds to the XXZ chain at ∆ = −1, while the state at α = 2 is the ground state of the Haldane-Shastry Hamiltonian [30,31]. The case α = 1 corresponds to a (symmetrized) Slater determinant, and its parent Hamiltonian is a free fermionic one (up to boundary contributions).
Participation spectrum
To obtain insights for generic values of α, it is instructive to rephrase Eq. (1) in the language of participation spectroscopy [58][59][60][61][62]. This consists of rewriting the wave functions Eq. (1) in a pseudo-energy fashion: In the last equality, we defined the function H α [{n}]: The functional coefficient E 0 [{n}] is an energy constant, while V (i, j) is a logarithmic interaction between occupied particles mediated by chord distances. Thus, we recognize H α [{n}] to be a 2D Coulomb gas (classical) Hamiltonian constrained in a 1D circular lattice [54,57]. Analogously, the wave function normalization Z α is a classical partition function: The parameter α plays the role of temperature and controls the leading weights in the JG states. The modulus squared coefficients in Eq. (1): are Boltzmann weights with classical Hamiltonian 2H α and partition function Z 2 α . The pseudo-energies of 2H α are collectively named participation spectrum and denoted ε({n}).
The ground state ε min determines the larger weights in the sum Eq. (1). For α > 0, the Hamiltonian favors repulsion among particles, constrained by the half-filling condition. Thus, the most probable configurations come from alternating occupation numbers. At negative temperature α < 0 the dominant coefficients are those maximizing the number of occupied nearest neighboring sites. For both cases, such configurations are not unique but degenerate, and for large values of α these states are expected to be the most relevant contributions to the Jastrow-Gutzwiller wave functions. Consequently the JG state are captured by the coherent superposition of these degenerate configurations, which leads, for α 1 and α −1, respectively to antiferromagnetic and ferromagnetic Greenberger-Horne-Zeilinger (GHZ) states [63]: The former state is usually dubbed Néel/anti-Néel state and corresponds to a global Schrödinger cat state. Apart from these extreme limit, at intermediate values of α the system exhibits competing weights, which render analytical arguments demanding if not impossible.
To test this heuristic argument, we consider the gap G = ε min − ε 1 st between the ground state energy of Eq. (4) and its first excited energy, which we refer to as participation gap. Let us discuss the case α > 0. It is convenient to introduce the number of ferromagnetic domain walls as the number of consecutive occupied/unoccupied sites N dws . For example N dws (|010101 ) = 0, while N dws (|011001 ) = 2. The Néel and anti-Néel states, i.e. the most The participation gap G increases linearly with α, with a coefficient that saturates to a constant g ∞ already at modest system sizes. (d) Pseudo-energy differences between two domain walls as a function of the domain-wall separation r. This is a measure of the confining potential between domain walls. The black solid line is the Luttinger liquid prediction [60] with Luttinger parameter K = 1/α. The fit describe extremely well our data, for 4 ≤ L ≤ 64.
probable states, are the only ones with N dws = 0, and all other pseudo-energy excitations can be easily labelled with this number. In Fig. 1 we present the participation spectrum of the JG states for α = 2, 6 and L = 16. The gap G between the most probable and the second most probable state increases linearly with α, with an exactly computable L-dependent constant g L . This saturates a thermodynamic value 1 g ∞ already for modest system sizes.
It is important to emphasize one aspect that is relevant in determining the system properties in the thermodynamic limit. The ground state pseudo-energy with alternating occupied sites is doubly degenerate for every system size. Instead, although the configurations with domain walls are exponentially suppressed in α, their degeneracy scales linearly with system size. In particular at L ∼ exp (cα) for some constant c, we expect a competing and non-trivial behavior between the Néel sector and the first excited sector. This has potentially relevant consequences, which are difficult to predict with the present study. In particular is unclear what effect this pseudo-energy thermodynamics have on quantum observables.
At a practical level, our results are consistent with the intuition above, that the Néel state predominately contributes for large α. In order clearly see the effects of the aforementioned thermodynamic competition for α = 6, we would have needed around L ∼ 10 4 sites. The large 1 We get an analytic expression for the constant: where the ellipsis indicate further computable digits.
gap for any computable finite L considered, renders these excited sectors negligible. The results for α < 0 are analogous to the latter, whereas the most probable configurations are the ferromagnetic ones and the excited pseudo-energy states are obtained as functions of antiferromagnetic domain walls, i.e. number of alternating occupied/unoccupied sites. However the most probable states there are L-degenerate: in the thermodynamics of the Coulomb gas this implies the low-lying pseudo-energy excitation are negligible even at small negative values of α.
Finally, from the substructure of the N dws = 2 sector we can extract how these domain walls interact. In particular, the pseudo-energy difference ∆ε 2dw = ε 2dw (r) − ε 2dw (2) between domains separated by a distance r and those close together (r = 2) has been used for local antiferromagnetic quantum Hamiltonian systems to distinguish between critical and symmetry broken phases of matter. In the former case, the domain walls are logarithmically confined with the separation distance; instead, in the latter this confining is linear. Moreover, the prefactor of this potential for 1D Luttinger liquids [66,67] is related to the Luttinger parameter. This has been tested in Ref. [60], where its authors analyze the XXZ chain.
Because of the explicit form of the classical Hamiltonian density Eq. (6), the interaction between two domain walls is expected to be logarithmic with their separation distance ( Fig. 1, panel (d)). By analogy with the XXZ phenomenology, one is tempted to conclude the JG states are gapless. If furthermore one assumes these states are representatives of Luttinger liquids, the fitted pre-factor suggests a Luttinger parameter K = 1/α. The latter statement has been recently conjectured [64]. This hypothesis is supported by CFT arguments [56] and from studies on the Resta polarization [65]. Here the authors estimate α c = 4 as critical value separating a conducting Luttinger phase to an insulating Néel ordered phase. Our data do not exhibits any transition point in the participation gap, nor a clear distinction between a gapped and a gapless phase. As remarked earlier, this may be due to a finite size effect, which we are not able to resolve at computationally affordable system sizes. In fact it is possible that the N dws = 2 domain walls sector results as decoupled for physical observables of the system, after a critical value of α. At present, however, the consequences of the participation spectroscopy to physical observables are unclear, and further studies are needed in this direction. In the next two subsections we improve our understanding of the Jastrow-Gutzwiller wave functions by numerically studying the entanglement entropy and the correlation functions of the Jastrow-Gutzwiller wave functions. We focus on these properties among others because they serve in the reconstruction technique and its quality checks. The considered system sizes suggests the existence of a critical phase between a Néel and a ferromagnetic GHZ regimes. Using finite size scaling we can bound the former in the interval α ∈ (0, 4.3).
Entanglement entropy
In this subsection, we discuss the entanglement entropy properties of the JG states (for related studies of Rényi entropies in 2D, see Ref. [14]). Entanglement is a fundamental quantity measuring quantum correlations among subregions of the system [68][69][70][71][72][73]. For pure states, this is determined by the spectrum of the reduced density matrix [74,75]. This operator is defined by giving a bipartition of the chain Λ = A ∪Ā and a state |Φ : Here ρ JG is the half-system reduced density matrix of the JG state. The results in green line are obtained through ED using symmetry restrictions. The red lines are the ferromagnetic GHZ predictions for the corresponding system sizes, while the black one is the Néel/anti-Néel cat state entanglement entropy.
Given its spectrum σ(ρ A ), we define the von Neumann entropy by: This function is a bona fide measure of entanglement for pure states when the Hilbert space factorizes in a tensor product form, H = H A ⊗ HĀ, and for this reason is usually referred to as entanglement entropy [76,77]. Fixing A = {1, 2, . . . , L/2}, we compute through exact diagonalization (ED) the von Neumann entropy for the state Eq. (1). We check the GHZ limits by comparing with the analytic calculations for the states in Eq. (9): The agreement is shown in Fig. 2. We isolate an intermediate region between the GHZ regimes by introducing the function: We plot this function in Fig. 3. Within this interval,S vN is logarithmic, with a pre-factor close to 1/3. This is consistent with exact solutions, where the systems display a critical regime. For instance, at α = 1 the system is a linear combination of Slater determinant. At this point the JG state correspond to a free fermion gas and the entanglement entropy can be computed analytically [78,79]: Here c is the central charge (c = 1 for free fermions) and the sub-leading term is a constant. The same scaling holds at α = 2, since the Haldane-Shastry Hamiltonian share the same universality class of the Heisenberg antiferromagnet [30,56]. By continuity, we argue the same critical behavior extends to the whole intermediate region. This is in line with the Luttinger liquid conjecture (see Sec. 2.2). Since the latter is of interest for the subsequent analysis, we estimate its bounding transition points. From Fig. 3 is clear that there is a transition in parameter space at α = 0.
We perform finite-size scaling on our data to estimate the critical value α c of the JG wave functions separating a critical phase with respect to a Néel ordered state. This is a phenomenological finite-size scaling procedure, since it is inherently related to a parameter characterizing the variational wave functions, and not associated to a coupling term in a Hamiltonian. Nevertheless, it is useful to bound the region of validity of the reconstruction method (Sec. 3), which relies on relativistic invariance. We consider the scaled entanglement entropyS(α) as an order parameter, as well as its derivative: which is roughly a susceptibility. We choose to consider both these quantities since the scaling we have is very mild with system size. From Fig. 3, introducing t = log(L) and α = (α − α c )/α c , we use the following simplified scaling ansatz: To perform the finite size scaling we vary the exponents ν, γ and the critical value α c over a suitable range of parameters. The fit is the best over different degrees of polynomials, test with a least-square method against the data [104]. By requiring the exponents to obey scaling relations γ = β − 1/ν we are able to reduce the fitting regime. We estimate the transition at α c = 4.3±0.1 with ν = 2.1±0.2 and β = −0.15±0.3. Value and error bars are the average and standard deviations of the best fits varying the range of system sizes considered. In Fig. 4 we plot both the order parameters of interest and the optimal data collapse. While the quality of the collapses is generically good, the modest system size are not able to resolve more efficiently the exponent landscape, which results quite flat. We believe a more systematic analysis is needed to better characterize the entanglement entropy and its phase transition for the JG wave functions. This would be a useful test also for the Luttinger liquid conjecture in Ref. [65], where it is argued the transition is around the value α conj c = 4. In this paper we choose to follow a more restrictive and cautious approach, focusing on subintervals of α ∈ (0, 4) in the rest of the paper. A concluding remark, which will be useful later, is about the α = 0 point. As previously discussed in the context of participation spectrum, this point is peculiar since the JG state is in an equal-weight combinatorial superposition. Its exact entanglement entropy can be computed [80]: We see that the pre-factor is different from the one in Eq. (15), signal that the state is not representative of the same phase. One can see this by investigating the properties of the exact parent Hamiltonian at α = 0: the XXZ chain at the ferromagnetic transition [56,80]. This Hamiltonian has a gapless quadratic spectrum, thus it breaks relativistic invariance due to a different dynamical exponent 2 z = 2. This observation will be important when trying to reconstruct local Hamiltonians using a relativistic ansatz. Indeed, as we shall comment in Section 4, for α = 0 the algorithm will not be able to return a correct parent Hamiltonian, as expected.
Correlation functions
To further characterize and resolve the Jastrow-Gutzwiller states, we compute the one-body and two-body spin correlation functions {σ z , σ + , σ − }. Their scaling properties resolve the nature of the state being critical or not. Due to the binary nature of the n i variables, for notational convenience we introduce the unary-not operator F ij acting on the site i, j, whose action on basis state is defined by logical negation on n i and n j . Since the system exhibits a U (1) symmetry related to number conservation, we compute only U (1) invariant correlation functions. Recalling σ z = 2n − 1 with n the number operator we have: At half-filling the first one is identically zero. The latter ones can be easily implemented numerically. The correlation length can be extrapolated through finite size scaling of the connected correlation function σ z i σ z i+L/2 c : Here a is a constant, while γ characterize the algebraic decay. In all the above equations, we exploited periodic boundary conditions. Let us stress that the definition Eq. (22) is meaningful only when the cluster decomposition principle holds. This requires the connected correlation function to decay to zero with the distance between the spins. This definition is used throughout in the literature of critical phenomena, where the phase is defined through the ground state manifold of specific Hamiltonians [82]. In this context, symmetry broken phases at finite system size manifest themselves as a coherent superposition of the ground states in the different symmetry sectors (GHZ states) [23]. The latter are a remarkable example of states which do not respect the cluster decomposition. To avoid odd/even effects, we present only L multiples of four.
Having the above remark in mind, we consider the definition Eq. (22) to also characterize the parameter space of the JG wave functions. Here we first check the system is fulfilling the cluster decomposition principle condition. When this is not the case, we expect the JG state to be representative of a finite size symmetry broken phase. Within this setting, if the parameter ξ is finite, the exponential behavior dominates on the algebraic one and the system is gapped, while if ξ → ∞ the system behaves as critical.
In Fig. 5, we show the results of our fitting procedure, plotting the inverse correlation length versus 1/L. For the chain lengths considered, the thermodynamic limit is difficult to estimate since at finite size the inverse correlation length 1/ξ L can be trusted upon the value 1/L. However, all values α < 4.0 are compatible with an infinite correlation length.
For large positive values and negative values of α, the cluster decomposition principle fails. The corresponding GHZ states (introduced in Sec. 2.2), representatives of symmetry broken phases, are confirmed to reproduce the correlation functions of the JG wave functions. A detailed discussion is given in Appendix A.
Entanglement guided search for parent Hamiltonians
In this section we summarize the scheme we employ to reconstruct parent Hamiltonians [50]. As previously remarked, this method requires additional conditions to work. This in contrast to other techniques [45,46] based on the quantum covariance matrix (QCM). The latter are simpler to implement since are based on requiring the input state to satisfy the zero energy variance condition. Thus, those methods generically guarantee that the input state is an eigenstate (not the ground state) of the parent Hamiltonian. Here comes the reason we have chosen to use the Bisognano-Wichmann Ansatz (BWA) scheme: the additional physical constraints guarantee the parent Hamiltonian of the input state as the ground state. This condition is at the core of eventual simulation protocols, since excited state are less robust in analogue experiments. Nevertheless, the relativistic requirement can be applied only to a narrow number of settings: for example if non-translational system are considered, such as disordered systems, BWA fails while QCM still gives meaningful results [103], provided a a fortiori analysis is done on the parent Hamiltonian space and their spectra. The method we adopt is based on the Bisognano Wichmann (BW) theorem, which for convenience we recap in the first subsection. Then, we introduce the common ingredients shared with other aformentioned techniques [45][46][47]49,50]. We conclude this section by presenting the algorithm and our chosen implementation.
Bisognano-Wichmann theorem and lattice models
By definition, reduced density matrices are positive operators with bounded spectrum σ(ρ A ) ⊂ [0, 1]. Consequently, it is always possible to find a lower bounded operator K A such that ρ A ∼ exp(−K A ). This object is usually referred to as entanglement or modular Hamiltonian, and in general is highly non-local, being the logarithm of the non-local operator ρ A .
Remarkably, Bisognano and Wichmann proved that the entanglement Hamiltonian acquire a local density when considering the ground state of a relativistic quantum field theory partitioned into two half-spaces [52,53,77,81]. Moreover, the density of this modular operator is proportional to the one of the theory Hamiltonian. The statement is the following.
Theorem (Bisognano Wichmann) Given a local relativistic QFT in d + 1 spacetime dimensions, described by an Hamiltonian H = d d xH(x) the half-space reduced density matrix of the vacuum |Ω is: Here A and B are respectively the manifolds A = {x ∈ R d : x 1 ≥ 0} and its complementary, while v is the sound velocity of the relativistic excitations. Sometimes, the pre-factor β ≡ 2π/v is dubbed entanglement temperature due to the analogy with respect to thermal density matrices.
More recently, this result has been revisited in the context of holography and many-body physics [83][84][85][86][87][88][89][90][91][92][93][94]. In particular, the theorem has been extended for theories with conformal invariance [83,85,86]. Given the subsystem A = {x ∈ R d |0 ≤ r ≤ R, r = ||x||}, its entanglement Hamiltonian reads: Interestingly, when considering lattice systems exhibiting relativistic low-lying excitations, the discretisation of Eq. (23) and Eq. (25) gives a fine approximation of their reduced density matrices [78,93,[95][96][97][98][99][100], with even exact results for specific models [101,102]. Moreover, the discrepancies due to the lattice structure disappear in the thermodynamic limit. This motivates the core idea behind the BWA method: to find optimal BW entanglement Hamiltonian describing the reduced density matrix of state of interest, in our case the Jastrow-Gutzwiller wave functions. For concreteness, in the remaining of this paper we make use of the discrete version of Eq. (25) in 1D system of size L and A = {1, 2, . . . , L/2}: Here r label the sites, h r is the lattice density of the Hamiltonian H, while K A the corresponding modular operator. Conventionally, we chose to absorb the entanglement temperature in the Hamiltonian density couplings h r .
Basis of local operators
To quantitatively describe the theory and entanglement Hamiltonians on the lattice we introduce the basis of local operators. As previously mentioned, these fully characterize the operator space of the parent Hamiltonian search. We say an operator is k-local if either (1) it has finite domain k-nearby few body operators, or (2) it is written as a linear combination of the latter. Furthermore, we require k to be constant for any finite system size L we consider. If these conditions are not fulfilled, we say the operator is non-local.
We define a basis of k-local operators as the set of matrices {O µ,r } µ∈I,r∈Γ . Here I is a set of internal indices, while Γ ⊂ Λ is a set of sub-lattice ones. Depending on the values of I and Γ, these basis span different vector spaces of local operators, whose generic element is: The dimension of these spaces is thus given by the combined cardinality of the label sets D = |I||Γ|. Before moving on, we clarify the above notation through few examples. Let us first consider the Pauli algebra at each site r ∈ Γ = Λ: The generic linear combination is: We see the total dimension is D = 4L in this case. A less trivial example is the two-body nearest neighboring interactions: Here α covers, in addition to the elements in Eq. (31), the following two-body operators at each site r: O 4,r = σ x r σ x r+1 , O 5,r = σ x r σ y r+1 , . . . , O 10,r = σ z r σ y r+1 , O 11,r = σ z r σ z r+1 .
The linear space has dimension D = 12L. Imposing symmetries one can reduce the dimension D of the operator space, in the same fashion symmetry constraints can be used to block diagonalize observables. For example, imposing U (1) and translational symmetry, a possible operator basis is the following: Here, the index α takes three values (D = 3) and the Hamiltonian is: In the second step of the above equation, we wrote the operators h α in terms of Eq. (31). Thus, the freedom of choosing the operator basis enables us to specify the required symmetries of the parent Hamiltonian, and it allows a reduction of complexity (for translational invariant systems, D ∼ O(1) in system size). Motivated by the symmetries of the JG states, we will consider the following basis for k ≥ 2: Varying the value of k we consider an increasing number of nearest-neighboring hopping and exchange operators. Finally, since the physics of the JG state at α = 2 is captured by a long range model, we shall consider the basis of non-local operators: These basis are both U (1) and translationally invariant, thus exhibits coefficients w α not depending on lattice sites. In literature, non-translational invariant basis have been employed in the reconstruction of disorder system Hamiltonians [45,48,103], or to enlarge the set of Hamiltonians having the input state as an eigenstate [46].
Parent Hamiltonian reconstruction method
We are now in position to present the BWA scheme. Let ρ input A be the half-system reduced density matrix of the the input state. We want to find optimal coefficients w α in Eq. (34) such that: This optimization can be implemented using any estimator of distance between ρ input A and the model reduced density matrix ρ BW A ({w α }). For example one can use the Kullback-Leibler divergence between the participation spectra of the reduced density matrices [60]. This estimator has the advantage of being easy to implement even for larger spacetime dimensions, but has the drawback of leading in general to a non-convex optimization. Such obstacle can be anyway surpassed using stochastic optimization algorithms. Instead, for the class of models described by the basis in Eq. (35) and Eq. (36), it can be proven that any convex estimator acting on the space of density matrices leads to a convex optimization problem (with a unique solution). Among these, we have found particularly useful for numerical implementations the relative entropy, which we adopt in the remaining of this paper. Given two density operators ρ and σ, it is defined as: This function quantifies the distance between between ρ and σ, it is non-negative S(ρ|σ) ≥ 0 (with the equality holding only if ρ = σ) and it is jointly convex. In particular, its restriction to a single argument is a convex function. As already stated, the relative entropy leads to a convex optimization admitting, up to numerical precision, a unique solution [50]: The relative entropy value express a "distance" in the reduced density matrix manifold, and quantify the difference between the initial wave function and the closer one fulfilling the BW theorem.
We implement a gradient descent on the relative entropy. Introducing the notation ∂ α = ∂/∂w α and: the gradient of the relative entropy reads We remark that the actual input needed are just the expectation values over the ground state and over the "thermal" BW density matrix. The former can be sometimes computed analytically, as in the JG states (see Section. 2), while the latter can be implemented with different numerical methods, including quantum Monte Carlo when no sign problem is present.
Reconstruction of Jastrow-Gutzwiller parent Hamiltonians
In this section, we apply the entanglement based reconstruction technique to JG wave functions, considering different choices for the operator basis. We quantify the quality of the reconstruction utilizing (1) relative entropies between reduced density matrices, (2) wave function overlaps, and (3) correlation functions. In view of the discussion in section Sec. 2, we focus here on the regime 0 < α < 4; the regimes where the wave functions are captured by GHZ states are instead discussed in Appendix A.
Models for reconstruction
We consider two paradigmatic classes of operators as candidates for the parent Hamiltonian reconstruction. The first one are the k-local Hamiltonians constructed from the basis B N N (k) These Hamiltonians for k ≤ 4 are archetypal for the study of strongly correlated matter in 1D and 2D, and have been used for ab initio numerical studies of quantum spin liquid phases in different lattices [18-21, 51, 104]. We notice that these operators contains the XXZ and the J 1 − J 2 model as particular cases. The second class are long-range XXZ Hamiltonians constructed from the basis B LR in Eq. (36): The reason in the latter choice is twofold: on one hand J 1 = ∆ 1 is the Haldane-Shastry Hamiltonian, the exact parent Hamiltonian at α = 2. On the other hand, in Ref. [54] Shastry conjectured that α = 2 is the ground state of Eq. (43). We remark that the parent Hamiltonian is defined up to an overall multiplicative constant which sets the energy scales, and an additive zero energy value. Thus, without loss of generality, we factor out the J 1 term and we are interested in the values {w/J 1 }.
Numerical implementation
We search parent Hamiltonians of the above form through the BWA technique. The implementation is based on exact diagonalization (ED) routines in Fortran, using standard libraries and LAPACK [105]. We performed gradient descents with Figure 7: Relative entropy of the JG reduced density matrix and the BW converged one. We see that enlarging the domain of the operator involved, the quality of the results increases. The line α = 1 corresponds to a free fermions gas.
various threshold error th = 10 −3 − 10 −6 . In the considered region, we notice no qualitative change in the observable behavior, although a smaller threshold error requires more steps in the gradient descent convergence. For convenience, we present the results only for th = 10 −4 . At this value, the observables are determined with a precision of around 0.1%. The initial value of the couplings is drawn by a uniform random distribution on the interval [−2, 2]. Here the spreading plays a minimal role: since the optimal solution is unique (see Section 3), the only ambiguity is numerical and due to the truncation to th . The resulting uncertainty is in the last sensible digit of the relative entropy and of the other observables, which we lift through averaging over 50 initial configurations. As argued in Sec. 2, in the thermodynamic limit the system should exhibit a critical regime in the region α ∈ (0, 4.30). However, for the modest values considered L ∈ {4, 6, . . . , 20}, we chose to focus on the subregion α ∈ (0, 4), where finite-size effects are less severe.
Diagnostics for reconstruction
Let us introduce the observables we use to access the quality of the parent Hamiltonian reconstruction. Firstly, we evaluate the relative entropy S(ρ jas |ρ BW ) between the converged BW reduced density matrix ρ BW and the exact JG one ρ jas . Since this function is a "distance" in the density matrix space, it quantifies how much the BW density matrix approximates the input state.
We then introduce the module of the overlap | ψ jas |ψ rec | between the JG wave function |ψ jas and the ground state of the reconstructed Hamiltonian: We stress that this quantity is meaningful only for finite size systems, since it decays to zero in the thermodynamic limit, for any arbitrary small difference between two vector states (in analogy with orthogonality catastrophe [106]). Finally, we compute the following quantity, a cumulative estimate of how much the correlation functions over the reconstructed state differ from the exact ones: Here the first term is the correlation function respectively on the ground state of the reconstructed parent Hamiltonian and on the JG state eq (20). The 1/ (L) factor renders this object non-extensive, which is desirable when comparing different system sizes. For convenience, we call this operator the cumulative correlation difference. Equipped with these tools, in the following subsections we separately present the analysis for the previously introduced basis Eq. (42) and Eq. (43). On the former, we first discuss overlaps and relative entropies for different basis choices, and finally discuss correlation functions. On the latter, we focus the analysis only on the relative entropy.
Reconstruction with N N (k)
We begin by considering the models in Eq. (42) for k = 2, 3, 4. If a p-local Hamiltonian exists, we expect the terms k > p to be finite size terms and to decay to zero enlarging the system size. We anticipate that our result suggests that an exact local parent Hamiltonian exists only for α = 1 (see, e.g., the scaling of the overlap depicted in Fig. 8), which corresponds to free fermions 2-local Hamiltonian. At different values of α, the reconstruction is only approximate, although it improves considerably increasing the basis N N (k). We deduce that the exact parent Hamiltonian should involve long-range interactions. Search for nearest-neighbor Hamiltonians. -Let us first restrict the easiest setting, that is choosing the N N (2) basis. In this case, the Hamiltonian Eq. (42) corresponds to the XXZ model. The value of interest is ∆ 1 /J 1 . When this is zero, the model reduces to the XX chain, which is a free fermion model up to a Jordan Wigner transformation. Moreover, it is interesting to compare our results with those of Ref. [56]. There, the authors considered the inverse variational problem, optimizing the parameter α with respect to the fixed ratio of ∆ 1 /J 1 . They argue that for α ∈ [0, 2] the wave functions are representatives of the critical phase ∆ 1 /J 1 ∈ [−1, 1] characterizing the spin-1/2 XXZ chain. Our results are compatible with their findings and the analytic results (Fig. 6).
For larger values of α, our results still indicate a very clear convergence to the thermodynamic limit. Moreover, the extrapolated values ( Table 1) always indicate that ∆ 1 > J 1 in this regime: this is compatible with an antiferromagnetic state with a very large correlation length. This finding is highly non-trivial, as there is no guarantee that our method shall return the correct parent Hamiltonian even in the presence of strong finite-volume effects, that have to be expected in this regime since, in the XXZ model, the transition to an antiferromagnetic phase belongs to the Berezinskii-Kosterlitz-Thouless universality class.
Search beyond nearest-neighbor Hamiltonians. -It is important to test the stability of these findings both with respect to enlarging the basis, considering N N (k > 2), and to system size. We thus considered the reconstruction also N N (3) and N N (4), and studied the behavior of the couplings {w α /J 1 }. As shown in Fig. 7 and Fig. 8, both the relative entropy and the overlap improve including higher-k terms. In addition, the magnitude of the couplings corresponding to the latter seems to increase with system size (see Fig. 9), suggesting that the exact Hamiltonians for the Jastrow-Gutzwiller states are long-ranged. An exception is the point α = 1, whose reconstructed Hamiltonian converges to the XX chain. As argued in Sec. 2, this is expected due to analytic arguments.
Ferromagnetic JG wave function. -Another particular point is α = 0. There, the corresponding JG wave function is the exact ground state of the ferromagnetic transition point XXZ. The BWA in principle should not work being this point described by a nonrelativistic field theory [80]. However, the converged coupling is flowing toward the correct ∆ 1 /J 1 = −1 enlarging the system size. Importantly, this result is strongly dependent on the basis chosen, and we see that it is unstable adding larger hopping terms (N N (3) and N N (4)).
Here the modulus of the couplings corresponding to (k > 2)-local terms increases, signal that a relativistic exact parent Hamiltonian for this point, if it exists, it is strongly long-range.
Correlation functions. -Finally, we present in Fig. 11 the results for the cumulative correlation difference V (rec|jas). At fixed system size L, it slightly increases when including higher k-terms. This is counterintuitive, since we observe that a larger basis N N (k) leads to states that are more similar to the JG wave functions (see Fig. 7 and Fig. 8). With the present analysis, we are not able to fully characterize if this trend is due to finite size effects or it has a more systematic nature. A possible explanation would be hidden in the BWA algorithm: since it optimizes over the short-k correlations (see Eq. (41)), the large distance correlators are less controlled and are subject to frustration effects. Within this interpretation, these discrepancies may suggest that longer range terms are required in the optimization to faithfully reconstruct an exact parent Hamiltonian.
Instead, at a fixed value of k, the cumulative correlation difference seems to saturate at some finite value. Being such an object deviation measure from a standard value (see Eq. (45)), it roughly gives how much on percentage the correlation functions change at a fixed site. In the worse scenario of our results, this has a value of around 10%. One may compare our findings with the exact results of the Haldane-Shastry model and the antiferromagnetic Heisenberg chain [30,31,107]: From the latter equations, we read the relative error of the nearest neighboring correlators and next-nearest neighboring ones, respectively of 2% and of 8%.
Combining the above reasonings, we state the reconstructed parent Hamiltonians are only approximate and the true parent Hamiltonians for the JG states require non-local terms. This further confirms our previous analysis. The exception is the point α = 1, where the cumulative correlation difference improves both with system size and by including larger N N (k). Figure 9: Ratios of the converged couplings ∆ 2 /J 1 and J 2 /J 1 versus inverse system size. We see that α = 1 is flowing toward the XX Hamiltonian (see also Fig. 6), while the other converged values are stationary in non-null values, suggesting long range 2-body physics for the JG states. Figure 10: Ratios of the converged couplings ratios versus inverse system size using the N N (4) basis. As for the previous cases, we see that α = 1 is flowing toward the XX Hamiltonian (see also Fig. 6,Fig. 9). Other values of α suggest a long range 2-body physics for the JG states. Relative error of the variational energy -As a last check we compute the variational energy of the parent Hamiltonian with respect to the Jastrow-Gutzwiller input state: and compare with the exact ground state energy E gs . The results are quantitatively compared via the relative error: We present the our results in figure Fig. 12. At fixed value of N N (k), our data suggest a mild linear growth of the relative error with system size. A linear extrapolation of the thermodynamic limit is given. All the considered cases lie within 1% of relative error in the energy landscape. Interestingly, at fixed L the relative error increases including larger N N (k), in a similar fashion to what we observe in the correlation functions. At present we cannot fully understand and characterize such counterintuitive behavior. As already mentioned in the previous paragraph, this may be due to the algorithm forcing the optimization on a finite size landscape and creating frustration effects. The latter likely explain the case α = 2, which should converge to the Haldane-Shastry pre-factors. Another possibility is that a new operator content is needed, and the chosen basis cannot grasp the thermodynamic properties of the systems. Further investigations on this problem are left for future studies.
Reconstruction with the long range model
We investigate the reconstruction when considering the model Hamiltonian Eq. (43), limiting our discussion to the relative entropy detector (see Sec. 4.2). The couplings are reported in Fig. 6, compared with the N N (k) cases. For the chain lengths considered, only at α = 2 the relative entropy shows a decreasing trend with system size (Fig. 13). This indeed corresponds to the exact Haldane-Shastry parent Hamiltonian. However, except at this fine-tuned point, At present we cannot infere a clear thermodynamic behavior for k → ∞, as the number of points is too modest to fit. Nonetheless, the error seems bounded within 1%-2%.
the relative entropy grows with system size, suggesting the parent Hamiltonian Eq. (43) is no the exact parent Hamiltonian for α = 0, and other more intricated terms must be added. Figure 13: Relative entropy between converged BW density matrix and the JG one for the long range model Eq. (43) for L = 8, 12, 16, 20. The results show a decreasing relative entropy for α = 2, which suggests the algorithm is approaching thermodynamic convergence. Instead, even points close to this Haldane-Shastry point exhibits increasing entropy, and certifying only an approximate reconstruction.
Conclusion and outlooks
In this work, we reconstructed approximate parent Hamiltonian for the one-dimensional Jastrow-Gutzwiller wave functions. We identified a region in parameter space where these wave functions display critical properties. Outside this interval, they are effectively described by Schrödinger cat states. Most likely, they are representatives of symmetry broken phases and their parent Hamiltonian is classical and constrained by the half-filling condition on the states. For the reconstruction technique, first we considered k-local Hamiltonians. We confirm the exact point α = 1 corresponding to free fermions, obtaining the XX Hamiltonian. At α = 0 the method fails to find local and relativistic parent Hamiltonians. This is due to a breakdown in the relativistic invariance in the wave function, whose exact parent Hamiltonian manifest gapless quadratic spectrum [56,80].
Our findings suggest the exact parent Hamiltonian for α = 1 should involve more complicated U(1)-invariant interactions, potentially with larger support. We checked the hypothesis of Shastry (Ref. [54]) of considering long-range XXZ chains with square secant couplings. Up to the considered system size there is a slow trend toward larger relative entropy, thus suggesting the ansatz is likely to be insufficient. Nevertheless, finite-size results are of value for Hamiltonian engineering and quantum simulations. Indeed, the BWA method provides inherently finite-size optimization and control on the basis chosen and on the quality of the outputs. In particular one can choose experimentally suitable operators in the basis, such as two-body operators. The fact that our technique is easily adaptable to include fully-longranged interactions may also be used in a different manner, that is, to certify and validate quantum simulators aimed at finding ground states of spin models including slowly-decaying power-law interactions, which are realized in both trapped ions [27] and Rydberg atom experiments [25,108].
It is of primary interest to apply similar techniques and considerations to two dimensional wave functions, such as the Laughlin wave functions. In fact, being the only computational demanding part of the algorithm the calculation of the ground state and the Bisognano-Wichmann expectation values, in principle one can tackle also higher dimensions by using Monte Carlo techniques. From the quantum engineering viewpoint, another intriguing perspective is to search for Liouvillians that have Jastrow-Gutzwiller wave functions as unique steady states [109,110]. In particular, dissipation may considerably soften the requirement for long-range couplings thanks to correlations induced by the bath.
A Correlation functions and parent Hamiltonian for the GHZ regimes
We argued that the JG states at α < 0 and α 1 corresponds to ferromagnetic and antiferromagnetic cat states. A first check is given by means of the participation spectrum and of the entanglement entropy (see Fig 2 in Section 2). Given the simple form of these GHZ states Eq. (9), we can compute their analytic correlation functions: In Fig. 14 we check the agreement between the above equations and the numeric correlation functions computed on the exact JG states. Our results suggest the state is in a symmetry broken phase [23]. Intuitively, we can guess classical parent Hamiltonians having these states as the ground state. For example, a ferro/antiferro-magnetic Ising model with the constraint of having zero magnetization. In practice, one can represent these states as MPS and use well-known results [11,23] to reconstruct local parent Hamiltonians. Figure 14: Difference between numerical correlation functions computed on the JG states and the analytic formulae Eq. (49). The different system sizes show a scaling to zero, confirming the correctness of the GHZ limit. | 11,779.6 | 2019-09-25T00:00:00.000 | [
"Physics"
] |
Gradient voltage amplification effect in FDSOI NCFET with thickness-variable ferroelectric layer
In this paper, a negative capacitance field effect transistor with thickness variable ferroelectric layer (TVFL NCFET) based on the fully depleted silicon on insulator (FDSOI) is proposed. The TVFL NCFET features the linearly increased ferroelectric layer thickness along the channel from source to drain. The gradient voltage amplification effect caused by the TVFL is analyzed according to the proposed capacitance model and simulation. Both of the model and numerical results indicate that the TVFL leads to a gradient increased electrostatic potential distribution along the bottom of the ferroelectric layer. The influences of gradient voltage amplification effect on the transfer characteristics, the output characteristic, the ratio between on-state-current (I ON) and off-state-current (I OFF), the drain induced barrier lowering (DIBL) and the subthreshold swing (SS) are investigated. The results show that the TVFL NCFET achieves the SS of 53.14 mV/dec, which is reduced by 19% when compared to the conventional NCFET. Meanwhile, large ION/IOFF is also realized and up to 1012 at most.
Introduction
The negative capacitance field effect transistor (NCFET) has been widely investigated in recent years because of its ability in improving subthreshold swing (SS) and reducing power dissipation.The operation of negative capacitance introduced by ferroelectric layer and other capacitances in the NCFETs is described by the capacitance matching [1,2].Meanwhile, the ferroelectric layer in NCFET is expected to realize the low SS [3].The ferroelectric materials doped hafnium oxides like Hf x Zr 1-x O 2 (HZO) are wildly used as ferroelectric layer to improve the performance and compatibility [4][5][6][7][8][9].The 2D thin film channel is also considered because it's excellent ability in achieving steep SS and large I ON /I OFF ratio [9][10][11].The special structure such as T-shaped gate NCFET [12], highly-doped double-pocket double gate NCFET [13,14] have been proposed and shown great ability in achieving steep SS [15] introduces a box ferro FDSOI which attaches ferroelectric layer to buried oxide to mitigate the negative differential resistance (NDR) effect.The first principle explanation of the NDR effect is presented, and the element which can influence the characteristic of NCFETs is also investigated in this reference.The thickness of the ferroelectric layer plays a role in the electric characteristic of the NCFET [16] indicates that the thicker ferroelectric layer brings steeper SS but larger hysteresis.
In this paper, a novel negative capacitance field effect transistor with thickness variable ferroelectric layer (TVFL NCFET) is proposed.The TVFL can result in a gradient increased electrostatic potential distribution along the bottom of the ferroelectric layer, named gradient voltage amplification effect.Which influences the device characteristics especial reduces the SS.A capacitance model is established to reveal the mechanism of the gradient voltage amplification effect caused by TVFL.The Sentaurus TCAD tool is used to investigate the gradient voltage amplification effect and its influences on the device characteristics.By calibrating with experiment baseline FDSOI in [17], the simulation method and models are determined, the models used in this work including the effective intrinsic density, doping-dependent mobility, doping-dependent SRH, Auger, high field saturation and FEPolarization.
Device structure and mechanism
The structure of the TVFL NCFET investigated in this paper is shown in figure 1.The TVFL NCFET characterized by a ferroelectric layer with linearly increased thickness along the channel from source to drain.Zirconium doped HfO 2 (Hf 0.5 Zr 0.5 O 2 ) is used as ferroelectric material.The Landau parameters are from [18], where α = −5.810× 10 10 cm F −1 , β = 3.286 × 10 19 cm 5 /(F • C 2 ) and γ = 2.165 × 10 28 cm 9 /(F • C 4 ).The thickness of the ferroelectric layer is denoted as T FS near the source and denoted as T FD near the drain.T FS = T FD means the conventional NCFET that has a uniform thickness of the ferroelectric layer.The TVFL NCFET is based on the fully depleted silicon on insulator (FDSOI) technology with the box thickness of 25 nm and top silicon layer thickness of 4 nm.The doping concentration in the source and drain is 1 × 10 20 cm −3 , and in the channel is 1 × 10 17 cm −3 .The gate length and gate oxide thickness are 20 nm and 0.6 nm, respectively.The overlaps between the gate and the source/drain region are the same value of 1 nm.
The ferroelectric layer is defined by LK model [19].The LK equation is expressed as: where ρ is the viscosity associated with polarization-switching dynamic.The free energy U of the ferroelectric layer is expressed as: In this equation, α, β and γ are ferroelectric material parameters, g is a coupling coefficient for polarization gradient term of the free energy.P is the polarization intensity.E stand for external voltage.Then, combining equations (1) and (2), the expression of E can get as followed: Under single-domain states condition, dP/dt = 0 and g = 0, so equation (3) could be simplified as followed: By giving ferroelectric parameters α, β and γ, the polarization characteristic of the ferroelectric layer could be decided.
The definition of electric field is E = V FE /T FE , where V FE is the voltage drop across the ferroelectric layer.Therefore, V FE can be expressed as: By ignoring the high order terms of P, equation (5) can be simplified as: The polarization intensity is numerically equal to the charge density generated by the polarization on the dielectric surface.In a unit area, the polarization intensity is numerically equal to the charge generated by polarization.So the P in equation ( 6) can be changed to Q [3].The capacitance of ferroelectric layer (C FE ) is defined by C FE = ∂Q FE /∂V FE .By the definition and equation (6), C FE can be approximatively calculated: According to the capacitance model shown in figure 2(a), the voltage amplification factor A V can be expressed as followed: where C mos stands for C OX , C Source and C Drain .
Combining equations ( 7) and (8), the voltage amplification A V is related to ferroelectric layer thickness T FE , express as: Normally, lager T FE induces greater voltage amplification.However, changing gate voltage leads to variation of depletion layer thickness at the top of the channel, resulting in variation in C dep in subthreshold region.In this investigation, the FDSOI technology makes a simplification which ignores the variation of C mos caused by changing T FD in this equation.
According to equation (9), the thicker T FE , the larger A V , which means the greater voltage amplification.For the TVFL NCFET, the ferroelectric layer thickness is linearly increased from source to drain.Therefore, the voltage along the bottom of the ferroelectric layer has been amplified and gradient increased from source to drain.Which is defined as gradient voltage amplification effect.To describe the gradient voltage amplification effect of the thickness-variable ferroelectric layer, a discrete capacitance model is proposed as shown in figure 2(b).In this model, the ferroelectric capacitance, oxide capacitance and depletion capacitance are separated into infinitesimal capacitance respectively.The larger ferroelectric thickness near drain side produces greater voltage amplification near drain.Therefore, V mos1 > V mos2 > V mos3 > K > V mos,n , and leading to V 1 > V 2 > V 3 > K > V n .These voltage differences create a lateral electric filed while it is caused by gate ferroelectric layer instead of drain voltage.They can bring extra current flow through infinitesimal resistance R 1 , R 2 , R 3 , K ,R n-1 in channel.This operation induces developed SS in subthreshold region and enhanced on-state current.
To verify the gradient voltage amplification effect, the Sentaurus TCAD tool is employed to simulate the TVFL NCFET. Figure 3 shows the electrostatic potential distribution in the ferroelectric layer.The T FS is fixed at 5 nm to clearly observe the negative capacitance effect, while the T FD is varied for 5 nm to 20 nm with equal step, standing for the conventional NCFET and TVFL NCFET with different linearly thickness.The voltage of 1.0 V is applied on the gate.As seen in the figures, the ferroelectric layers show a great ability of voltage amplification, which enlarge the gate voltage significantly.
Figure 3(a) is the electrostatic potential distribution in the ferroelectric layer of the conventional NCFET with T FS = T FD = 5 nm.It reached maximum voltage of 1.13 V at the bottom of the ferroelectric layer.Figures 3(b)-(d) are the electrostatic potential distribution in the thickness variable ferroelectric layer.It is obviously that the TVFL leads to a gradient increased electrostatic potential distribution along the bottom of the ferroelectric layer.The TVFL in figures 3(b)-(d) obtain the maximum voltage of 1.332 V, 1.772 V and 2.647 V at the bottom ferroelectric layer near drain side, respectively.As T FD increases, electrostatic potential gets higher, which indicates that the voltage amplification is stronger, and the gate control ability is better.
Figure 4 shows the electrostatic potential at the bottom of the ferroelectric layer for the TVFL NCFET with different T FD and fixed T FS .It's obviously that the voltage amplification near drain is larger than that near source, and the voltage gradient from source to drain is enlarged with the T FD increasing.It is worth noting that the potential at the source side is also improved with the increasing T FD although the T FS is fixed at 5 nm.This is because the ferroelectric layer thickness varies continuedly.From figure 2(b), it's clear that electrostatic potential on the top surface of channel V 1 , V 2 , V 3 , K, V n is determined by the complex coupling of each capacitance column, indicating that the electrostatic potential on the bottom surface V mos1 , V mos2 , V mos3 , K, V mos,n is relevant.
Results and discussions
The capacitance model and simulation results show that the thickness variable ferroelectric layer results in a gradient voltage amplification on the surface of the channel.In order to reveal the influences of gradient voltage amplification effect on the electrical characteristics of the TVFL NCFET, the output characteristics, transfer characteristics, DIBL effect, I ON /I OFF and sub-threshold swing are discussed.
Figure 5 gives the output characteristics of TVFL NCFET.The V GS varies from 0.3 V to 1.5 V with step of 0.3 V.In figure 5(a), it is found that the TVFL NCFETs induce a steep current increment as V DS gets larger at the liner region.But once the device enters saturation region, the current increment becomes slower and followed by current drop.This effect called negative differential resistance (NDR), gets more significant as T FD increases.This is because the voltage on the bottom surface of ferroelectric layer is influenced by both V GS and V DS through gate-drain coupling.The negative capacitance of ferroelectric layer introduces the negative correlation between V mos and V DS .The increasing V DS reduce the charges under ferroelectric layer and lead to increase of V FE in negative capacitance region, thus result in the reduction of V mos .This means the voltage amplification is impeded, which result in the decrease of I D .Meanwhile, according to equation (5), the ferroelectric layer thickness can be regard as a coefficient which modulate the change in charges.As T FD increases, the effect brought by the decrease in charges is amplified.In this case, the increasing T FD results in more significant NDR. Figure 6 shows the DIBL effect of TVFL NCFETs depends on changing T FD combined with the conduction band energy diagram.In the energy diagram, the V GS is set at 0.1 V to ensure the device is turned off.The V DS varies from 0.05 V to 0.3 V.It can be found that as V DS increases from 0.05 V to 0.3 V, the potential barrier in the channel region increases.This is a strong evident of reversed DIBL.Meanwhile, compared with the conventional ferroelectric layer, the TVFL brings a higher potential barrier, and the barrier increases more with the increase of V DS .As T FD gets thicker, the raise of conduction band barrier gets larger, proving that TVFL with larger T FD has greater ability in reversing DIBL and inhibit short channel effect.The reversed DIBL ensures the gate controllability at small V GS , which helps to reduce drain leakage current as shown in figure 5(b).It also can be seen in figure 6 that the DIBL is reversed to negative.As T FD increases, the negative DIBL become larger in absolute value.This result confirms the finding in conduction band energy diagram.
Figure 7 shows the on current and off current of TVFL NCFETs.It can be seen that as T FD increases, the on current increases, which means better current driving ability.This is because stronger voltage amplification induced by thicker ferroelectric layer contributes to the forming of conductive channel.It also can be seen that the off current decreases significantly as T FD increasing.It means that the leakage drain current is better restrained, which helps to reduce static power consumption.This can be explained by the reversed DIBL. Figure 8 shows the transfer characteristics of TVFL NCFETs.The logarithmic axis in I D is set.It can be seen that TVFL NCFETs achieve steep I D -V GS curves.With fixed T FS , the larger T FD corresponds to steeper I D -V GS curve, which means smaller subthreshold swing.In the subthreshold region, the larger gate voltage caused by the larger T FD contributes to the formation of the inversion layer, which promotes the increase of drain current during the transistor opening.In addition, larger T FD forms larger lateral potential difference under the gate stack and on the top of the channel, which also promotes the drain current.Thus, the steeper curve can be obtained and leading to the smaller subthreshold swings.Figure 9 gives the SS depends on T FD with fixed T FS .It's obviously that in TVFL NCFETs, the structure with larger T FD has better performance on reducing SS.This is because larger T FD induces greater voltage amplification on the drain side as explained in equation (9) and realizes better gate control and the lateral electric field which can induce extra current flow.The conventional NCFET with T FS = T FD = 5 nm achieves the SS of 65.56 mV/dec in forward sweep and 65.58 mV/dec in reverse sweep, while the TVFL NCFET with T FS = 5 nm and T FD = 20 nm achieves the SS of 53.14 mV/dec in forward sweep and 53.13 mV/dec in reverse sweep, which provides improvements of 12.42 mV/dec forward and 12.45 mV/dec reverse.
It seems like that thicker T FD means better performance, nothing different from increasing the thickness of whole ferroelectric layer.But actually, something is found when take consideration into the structure with T FS = T FD = 20 nm. Figure 10 22 mV maximum and nearly 0 mV at I D = 10 −7 A, also is negligible.For the NCFET with T FS = T FD = 20 nm, it has hysteresis of 237 mV maximum and 60 mV at I D = 10 −7 A, which brings about severe problem of hysteresis and will lead to serious logic confusion in the circuits.In contrast, the TVFL NCFET with T FS = 5 nm, T FD = 20 nm shows hysteresis by 91% and the hysteresis happens at a lager voltage and away from threshold, which means that the impact to switching operation of the device is likely slighter.Therefore, the thickness of the ferroelectric layer should be optimized and the TVFL is a useful method to restrain hysteresis while reducing the SS.
Conclusion
In this paper, the TVFL NCFET is proposed and investigated by theoretical model and TCAD simulations.It is found that the TVFL forms the gradient electrostatic potential distribution and has great ability on voltage amplification.Compared to the conventional NCFETs with T FE = 5 nm, the maximum electrostatic potential on the bottom surface of ferroelectric layer in TVFL NCFETs is increased by 98.5% at most.The TVFL NCFETs can achieve I ON /I OFF ratio of 2.12 × 10 12 , which is over 10 2 larger than conventional NCFETs with T FE = 5 nm.The SS of TVFL NCFETs is also reduced by 19% at most.Meanwhile, the simulation results show that compared to the conventional NCFETs with T FE = 20 nm, the TVFL NCFETs reduce hysteresis by 91% at least and can achieve non-hysteresis operation with larger performance loss in SS.
Figure 2 .
Figure 2. The capacitance model of NCFETs.(a) is the conventional capacitance model of NCFETs, which regards ferroelectric capacitance as a whole.(b) is the novel capacitance model of the proposed TVFL NCFETs, which separates ferroelectric capacitance into infinitesimal capacitance.
Figure 4 .
Figure 4.The electrostatic potential at the bottom of the ferroelectric layer.
Figure 5 (
b) zoom in the region of V GS = 0.3 V in figure5(a).The NCFETs are turned off at this V GS .It's obviously that larger T FD results in smaller I D .It means that the TVFL structure with thicker T FD can prevent device from punch through at small V GS .This phenomenon is reflected as negative (or reversed) DIBL.
Figure 5 .
Figure 5.The output characteristics of the TVFL NCFET that (a) V GS = 0.3 V, 0.6 V, 0.9 V, 1.2 V, 1.5 V respectively.(b) zoom in the region of V GS = 0.3 V.
Figure 6 .
Figure 6.The negative DIBL effect in TVFL NCFET.The conduction band energy diagram indicates the rising barrier when V G = 0.1 V and the device is turned off.
Figure 7
Figure6shows the DIBL effect of TVFL NCFETs depends on changing T FD combined with the conduction band energy diagram.In the energy diagram, the V GS is set at 0.1 V to ensure the device is turned off.The V DS varies from 0.05 V to 0.3 V.It can be found that as V DS increases from 0.05 V to 0.3 V, the potential barrier in the channel region increases.This is a strong evident of reversed DIBL.Meanwhile, compared with the conventional ferroelectric layer, the TVFL brings a higher potential barrier, and the barrier increases more with the increase of V DS .As T FD gets thicker, the raise of conduction band barrier gets larger, proving that TVFL with larger T FD has greater ability in reversing DIBL and inhibit short channel effect.The reversed DIBL ensures the gate controllability at small V GS , which helps to reduce drain leakage current as shown in figure5(b).It also can be seen in figure6that the DIBL is reversed to negative.As T FD increases, the negative DIBL become larger in absolute value.This result confirms the finding in conduction band energy diagram.Figure7shows the on current and off current of TVFL NCFETs.It can be seen that as T FD increases, the on current increases, which means better current driving ability.This is because stronger voltage amplification induced by thicker ferroelectric layer contributes to the forming of conductive channel.It also can be seen that the off current decreases significantly as T FD increasing.It means that the leakage drain current is better restrained, which helps to reduce static power consumption.This can be explained by the reversed DIBL. Figure 7 also gives the I ON /I OFF ratio of TVFL NCFETs with different T FD and fixed T FS .The logarithmic axis in I ON /I OFF ratio is set.The TVFL NCFETs show great ability in increasing I ON /I OFF ratio.With the same T FS , I ON /I OFF ratio increases as T FD gets larger.These results indicate that by choosing a proper V DS , the property of drain current can be improved greatly by the TVFL NCFETs.
Figure 8 .
Figure 8.The transfer characteristics of TVFL NCFETs with different T FD and fixed T FS .
Figure 7 .
Figure 7.The I ON and I OFF of TVFL NCFETs with different T FD and fixed T FS .
Figure8shows the transfer characteristics of TVFL NCFETs.The logarithmic axis in I D is set.It can be seen that TVFL NCFETs achieve steep I D -V GS curves.With fixed T FS , the larger T FD corresponds to steeper I D -V GS curve, which means smaller subthreshold swing.In the subthreshold region, the larger gate voltage caused by the larger T FD contributes to the formation of the inversion layer, which promotes the increase of drain current during the transistor opening.In addition, larger T FD forms larger lateral potential difference under the gate stack and on the top of the channel, which also promotes the drain current.Thus, the steeper curve can be obtained and leading to the smaller subthreshold swings.Figure9gives the SS depends on T FD with fixed T FS .It's obviously that in TVFL NCFETs, the structure with larger T FD has better performance on reducing SS.This is because larger T FD induces greater voltage amplification on the drain side as explained in equation (9) and realizes better gate control and the lateral electric field which can induce extra current flow.The conventional NCFET with T FS = T FD = 5 nm achieves the SS of 65.56 mV/dec in forward sweep and 65.58 mV/dec in reverse sweep, while the TVFL NCFET with T FS = 5 nm and T FD = 20 nm achieves the SS of 53.14 mV/dec in forward sweep and 53.13 mV/dec in reverse sweep, which provides improvements of 12.42 mV/dec forward and 12.45 mV/dec reverse.It seems like that thicker T FD means better performance, nothing different from increasing the thickness of whole ferroelectric layer.But actually, something is found when take consideration into the structure with T FS = T FD = 20 nm.Figure10gives the transport characteristics of NCFETs with different T FS and T FD = 5 nm.Both forward sweeping and reverse sweeping are done.The NCFET with T FS = 5 nm, T FD = 5 nm shows almost no hysteresis, but poor performances.The TVFL NCFET with T FS = 5 nm, T FD = 20 nm shows hysteresis of
Figure 10 .
Figure 10.The transfer characteristics of NCFETs with T FS = T FD = 5 nm, T FS = 5 nm, T FD = 20 nm and T FS = T FD = 20 nm.Both sweeping forward and reversely is included.
Figure 9 .
Figure 9.The SS of TVFL NCFETs with different T FD and fixed T FS in both forward and reverse sweeping condition. | 5,194 | 2024-04-17T00:00:00.000 | [
"Engineering",
"Physics"
] |
REALTIME BLOOD BANK DATA REPOSITORY
The real motive behind the Realtime Blood Bank Data Repository is to streamline and mechanize the way toward looking for blood if there should be an occurrence of crisis andkeepup the records of blood givers, bene iciaries, blood gift projects and blood stocks in the bank. As of now, people, in general, can just think about the blood donation event occasions through traditional media means like radio, paper or TV commercials. There is no data with respect to the blood gift programs accessible on any of the gateway. With the manual archive, there are issues in dealing with the records. There is no concentrated data set of volunteer benefactors. Along these lines, it turns out to be truly dreary for an individual to look through blood if there should arise an occurrence of crisis. This task plans to mechanize the blood and contribute to the board framework in a blood donation centre to improve the record executives effectiveness because of the developed size of records of information.
INTRODUCTION
The essential structural aim is to give blood donation administration to the city as soon as possible. Real-time blood donation centre data repository is an application that is Web-based and used to store, recover, measure and break down data about the authoritative and then stock up the executives of a blood donation centre.This venture targets keeping up all the data relating to blood givers, distinctive blood bunches accessible in each blood donation centre and assisting them with overseeing in a superior manner. Venture Aim is to give straightforwardness in this ield, make the interaction of acquiring blood from a blood donation centre problem-free and debasement free and make the arrangement of blood donation centre administration powerful. To do this we require a superior grade. Web Application to deal with those positions. Blood donation centre gift frameworks can gather blood from numerous donors in short from different sources and disperse that blood to penniless individuals who require blood. This is an electronic data set application framework that can be utilized by blood donation centres or blood focuses as a way to promote the cross country blood gift occasions to the general population and at the same time permit the medical clinics to ask for the blood. The framework keeps the record of the relative multitude of benefactors, bene iciaries, blood gift programs, and dismissed blood types.
LITERATURE SURVEY
Most blood donation centres are as yet running manual frameworks in their cycles. Thus, there is an absence of effectiveness since it is still paper-situated in gathering data about givers, inventories of blood sacks, and blood bonding administrations. The absence of appropriate documentation may jeopardize patients' well being because of the chance of polluting blood packs. Pollution happens when there is a de icient contributors' clinical history record and the blood sacks' time span of usability isn't observed as expected. Subsequently, an electronic blood donation centre administration framework may be expected to resolve these issues and issues experienced to guarantee blood bonding wellbeing.
Numerous examinations were held and led to the idea of blood donation centre, many of them worried about overseeing the records of benefactors to work with the cycles of gift, the others associated blood donation centres to one another in one framework and one information base, the rest utilized current innovation like an electronic card and standardized identi ication framework. Blood donation centre Information Management System is a data framework which assists with dealing with the benefactors records and patients at the blood donation centre.
It is mainly intended to recover, measure, store and dissect data which is worried about the authoritative and administration of stock inside a blood donation centre. Such sort of framework will permit the approved blood donation centre of icial to log in utilizing a mysterious secret phrase and effectively deal with the records of the blood donors and the patients who needs blood. Furthermore, the blood donation centre data the board framework isn't out of date to the experts; rather it assumes an extraordinary part in drawing in the contributors and different partners because of its effortlessness in the booking, and warning of gift time to the needy. A report published by Dr Sharad Maheshwari in the International Journal of Engineering Research and Application (IJERA), said that in India, the blood donation centre administration data framework MIS is a coordinated blood mechanization framework.
The electronic component which connects all the blood donation centres of each state into a solitary organization, The blood donation centre alludes to all obtaining, approval, capacity and course of different live information and data electrically with respect to blood gift and bonding administration. The ZNBTS aka Zambia Blood Transfusion Service with help from the International Institution for Communication and Development (IICD) has brought up an electronic framework that has digitized enlisting of the contributor and sends an SMS message to benefactors of blood advising them that they can give blood Again. In addition to that the product also makes it simpler to arrive at blood donors by registering their data and then saving it online in an information base to be opened from any of ice of the ZNBTS.
PROBLEM STATEMENT
The present-day situation in India is a paper-based system that is enduring lack of central data references, which results in a time-consuming process for retrieving data in addition to lack of security of data and human error which needs an alarming system to circumvent. One of the problems currently in hospitals which is very serious is the lack of blood supply during an emergency situation. The crucial need to transfer blood requires proper management to determine which blood group is available. The second problem which arises and that has led to the design of this system which in turn leads to human errors is that the vital information about the blood grouping, donor availability, tracing the data is complicated and time-consuming when it is done manually. Moreover, this system requires a lot of manpower, data is also not secure, in addition to that retrieval of data and production of reports are timeconsuming. In the past few years, India has seen an increment in the blood collection (from 9.5 million units in 2012 to 11.4 in 2017) . But still there was a shortage of about 1.95 million units of blood last year.
EXISTING SYSTEM
Various studies have been composed on the idea of blood donation centre administration frameworks with most of them commending computerization as a component of accomplishing effectiveness also, adequacy in this space accordingly not taking a gander at certain issues the framework may look due to restricted or abuse of functionalities. Pah Essah and Said Ab Rahman (2011) proposed an improvement of an administration data framework to oversee blood donation centre dependent on data of benefactor, bene iciary and blood. Their framework has three modules: the benefactor module, patient module and blood module.
Anyway, some essential issues are left to the side in this methodology, for example, who is liable for the organization of the framework. As indicated by Mailtrey D Gaijjart (2002) proposes an improvement of blood donation centre information framework as an answer for forestall close to miss occasions and improve record recovery. They contend that with computerization quick recovery of records will improve the productivity of blood donation centres activities. Akshay V Jain Khanter (2009) recommends an administration data framework application that covers a portion of the blood donation centre administration issues identi ied with a speci ic area. An intriguing methodology by Jeroen Benien and Hein Force (2012) is that of production network the executives for blood and blood items naming the cycle as sporadic and the interest for blood stochastic. This is of extraordinary rami ications if the administration of blood donation centres were to get compelling. At last, E. M. S. S. Ekanayaka and C. Wimaladharma (2015) fostered a blood bank's framework to accumulate all the blood givers into one spot naturally and illuminate them continually about the chances to give blood by means of an SMS to the benefactor's versatile telephone.
The following is a proposed framework that will dispose of the relative multitude of issues that the blood donation centres and board framework is confronting at present.
METHODOLOGY
The research methodology was built up by studying the problem de inition, gathering data, design, and inally conclusion. The process of the blood bank requires registration of donor, information of the donor and testing, donation of blood, bloodstock management, blood screening and stock movements, and presently, there's a manual system to keep count of the records of blood screening and donor registration which features a lot of problems and errors.
Also, it has tried to gather some forms that they use in collecting information of donors. The collected data is evaluated and analyzed and then the needed information is extracted.
DESIGN AND IMPLEMENTATION
As stated in methodology, the system proposed was divided into different screens, which represent the major departments in the blood bank. The system is used to maintain the full information about blood, donor and transfusion process in the laboratories, the system uses a database which is central to store all input details including blood stock, donor information and blood group's information. The proposed UML class diagrams are:
RESULTS AND ANALYSIS
In this section we show the results of our Realtime Blood Bank Data Repository model. Based on the analysis of information, database creation tables, division of modules, code the system results in the following screen's layout.
DISCUSSION
For the transition from paper based to electronic health record to be successful requires careful coordination, from implementation and selection to maintenance and training. A computerized Realtime blood bank data repository system is obtained after various steps which resulted in a variety of speci ications, shown as screens. The privacy screen is a major speci ication that facilitates users to save their entered information. Flexibility of system is an important facility which helps in adding more users and bene it when using Computerized Realtime Blood Bank Data Repository System, this also helps in sharing of different subunits. The system creates an integrated information domain where all the data of donor, compatibility of blood groups is available. As the world is directing towards new technologies, in a similar way this computerized system guarantees and ensures that every information entered to the system is valid to avoid errors. From the functionality point of view, the system ensures that less time is consumed, and from an economic point of view it minimizes the manpower, resulting in less expenses.
CONCLUSION
All in all, Computerized Realtime Blood Bank Data Repository System (CRBBDRS) may be a system that wants to control and manage all activities in bank departments.
The system records nad keeps all donor information, distribution of blood to hospitals, blood information etc. The implementation of the system was carried out in many steps. The system was designed by creating a database and connecting it with the programming language code which ends up in the system screens layout. The system also leads to meeting user requirements. In general, this project designs only a prototype of the Computerized Realtime Blood Bank Data Repository System (CRBB-DRS). This system has lexibility to modify itself and satisfy all needs to reach and cover up other departments as well, that is not currently covered here. | 2,641.8 | 2021-07-03T00:00:00.000 | [
"Medicine",
"Computer Science"
] |
Using Ethereum Smart Contracts to Store and Share COVID-19 Patient Data
Introduction The emergence and rapid spread of the coronavirus disease 2019 (COVID-19) pandemic have revealed the limitations in current healthcare systems to handle patient records securely and transparently, and novel protocols are required to address these shortcomings. An attractive option is the use of Ethereum smart contracts to secure the storage of medical records and concomitant data logs. Ethereum is an open-source platform that can be used to construct smart contracts, which are collections of code that allow transactions under certain parameters and are self-executable. Methods The present study developed a proof-of-concept smart contract that stores COVID-19 patient data such as the patient identifier (ID), variant, chest CT grade, and significant comorbidities. A sample, fictitious patient data for the purpose of testing was configured to a private network. A smart contract was created in the Ethereum state and tested by measuring the time to insert and query patient data. Results Testing with a private, Proof of Authority (PoA) network required only 191 milliseconds and 890 MB of memory per insertion to insert 50 records while inserting 350 records required 674 milliseconds and similar memory per insertion, as memory per insertion was nearly constant with the increasing number of records inserted. Retrieving required 912 MB for a query involving all three fields and no wildcards in a 350-record database. Only 883 MB was needed to procure a similar observation from a 50-record database. Conclusion This study exemplifies the use of smart contracts for efficient retrieval/insertion of COVID-19 patient data and provides a case use of secure and efficient data logging for sensitive COVID-19 data.
Introduction
The ongoing coronavirus disease 2019 (COVID- 19) pandemic has revealed the limitations in current healthcare systems to handle patient records. Specifically, the reliance on several entities to communicate sensitive patient information has raised questions related to security [1][2][3] and transparency [4][5][6][7], which has become even more pronounced during the pandemic [8][9]. Medical records of patients contain private information, such as past medical history, significant comorbidities, and lung CT scan grading. Access to such data not only allows for efficient medical treatment but also for health officials to track the progress of COVID-19. However, these abilities depend on the security of such records. Furthermore, as data of such nature are often held in databases using multi-level permission security, access and upload of time-sensitive information are affected. Blockchain technology is gaining traction in healthcare applications and has the potential to mitigate these issues.
Ethereum is a unique type of blockchain as it allows blockchain technology to adapt to software applications. For example, one could construct programs that run if and only certain requirements and conditions are satisfied. These programs are referred to as smart contracts, which are self-executable pieces of code that remain within the Ethereum state and allow certain transactions when called for [10]. Ethereum platform allows such code to be run by providing certain programmable logic requirements . The platform 1 2 2 2 3, 4 3 2, 5 6 7 also allows for data storage not limited to cryptocurrency. These characteristics allow for the platform and its smart contracts to be molded to other use cases that require robust integrity and amenability, such as holding patient data. For example, Ethereum blockchain has been shown to handle pharmacogenomics data comprising gene-variant-drug combinations, as well as neuroimaging data [11][12]. No studies, however, have so far investigated the implementation of Ethereum blockchain using COVID-19 patient data. The present study describes a case use of a smart contract to incorporate COVID-19 patient data into the Ethereum platform to enhance patient privacy, records security, and data immutability.
Materials And Methods
We present here a quick, efficient, and memory-saving smart contract to store and query COVID-19-related patient data in Ethereum. A brief overview of Ethereum is provided here. The Ethereum yellow paper contains additional technical explanations [13].
Explanation of smart contracts and the Ethereum system
The major data structure that Ethereum uses is modified Merkle Patricia tries. Instead of needing to store all the data in the blockchain, these tries allow the blocks to only store the root node hash of each trie, all the while preserving immutability ( Figure 1A). The major trie types in Ethereum include account storage trie, world state trie, transaction receipt trie, and transaction trie [13]. The world state is analogous to the global state that is frequently updated by executing transactions and maps account states and addresses (accounts). Information associated with certain accounts and smart contract data is in the account storage trie.
FIGURE 1: (A) Overview of Ethereum blockchain and (B) proof-ofconcept smart contract architecture
Smart contracts can be described as self-executable, Turing-complete: programs that perform within the Ethereum Virtual Machine (EVM) and maintain their own unique storage [13]. Operations within these contracts are allowed three space types for storing data. These types include memory, stack, and long-term storage. The contract's storage continues to exist after the computation has finished, as opposed to stack and memory [10]. Ethereum's native cryptocurrency, Ether, is used to pay costs; colloquially termed 'gas', these costs are associated with calling, storing, and retrieving smart contracts [10]. Smart contracts are programmed using an object-oriented programming language designed specifically for EVM, called Solidity [14].
Network nodes must agree as to which transactions are allowed to be added to the blockchain through consensus mechanisms. One such mechanism, Proof of Work (PoW) consensus, requires "mining" of blocks by nodes to solve a mathematical puzzle utilizing computational energy. Only when the puzzle is solved and validated by other nodes on the network is the block aggregated to the chain, and a reward is given to the node that succeeded in mining the block. A faster alternative to PoW is the Proof of Authority (PoA) consensus. This mechanism consists of certain trusted nodes, known as 'authorities', which are the only nodes allowed to add new blocks after validating transactions. Since this type of mechanism is faster than PoW, it is commonly utilized in configuring private blockchain networks [15].
Network configuration
Querying and storing sample patient data were experimented via smart contracts in Truffle suite v5.1.62, an Ethereum development framework equipped with a JavaScript testing environment. The development network was private and separate from the public Ethereum network, which allowed for testing without the need for constantly deploying. The contracts were run using the PoA consensus with six nodes on the network. As node configurations do not change performance in private PoA networks, different configurations were not assessed [16]. Default gas limits were utilized throughout smart contract testing. The smart contract code is available at https://github.com/Batchu-Sai/covid19.
Database architecture and insertion
Using arrays to store data requires iteration for observation retrieval. To circumvent this, mappings were implemented, which are conducive to efficient key-value lookups. Four storage mappings were implemented with Solidity's mapping type. Every observation was placed in its individual struct, which is a special type to group numerous variables of heterogeneous types. Fictitious patient data was used for one mapping. The other three mappings -the variant, comorbidities, and patient identifier (ID) -acted as keys to an array of counter IDs. This counter ID variable assigns a custom index for each observation insertion and is globally updated. This architecture allows using any of the three keys to retrieve the inserted patient observation using the specific counter IDs ( Figure 1B) and likewise allows for non-iterative lookup for more efficient database querying.
A function was written to insert a record for a single patient for a specific variant-patient ID-comorbidity combination ( Figure 2). Fictitious patient data was constructed, which included the COVID-19 variant, chest CT severity grade, patient ID, patient age, and significant comorbidities. If required, the function first converts the patient data to the desired storage type. Then the function checks if the combination already exists. If not, the variant-patient ID-comorbidity combination is aggregated to an array. The variant, patient ID, and comorbidities are utilized to key into their mappings. The counter ID is appended to the value array.
The counter ID variable is updated after the observation struct is inserted into the database with the keyvalue pair (comprised of the counter ID -patient data struct).
FIGURE 2: Pseudocode for inserting observation
The variant, CT scan grade, patient ID, and age were stored as bytes32 variables. The comorbidities field was treated as a string since comorbidities may be long descriptions and patients uncommonly present with several comorbidities.
Querying database
A function used to query the observations was written to take up to three fields (variant, patient ID, and comorbidities) along with wildcards to query by ( Figure 3). Initially, the function examines which fields have been queried. These fields are subsequently utilized as keys to retrieve the corresponding values from the counter ID mapping. If wildcards were used to query, then all the counter IDs are returned. For each record, the minimum length counter ID array is looped through to examine if it equals the counter ID from the outer loop. The matching ID is taken to collect the value of the struct from the corresponding database mapping. The framework was adapted from Gürsoy et al. [11].
Results
JavaScript console with external scripts tested the smart contract. To confirm accuracy, 50 stochastic queries were manually checked with the fields used to query and were confirmed to retrieve accurate results in every case. Time and memory for inserting and querying were also measured, with varying numbers of records stored in the database (Figure 4). The time to insert 50 records was 191 milliseconds while inserting 350 records required 674 milliseconds. Inserting 50 records required 890 MB of memory per insertion. To insert 350 records required nearly the equivalent memory per insertion ( Figure 4A). Retrieving required nearly 912 MB for a query involving all three fields and no wildcards in a 350-record database. Only 883 MB was needed to procure a similar observation from a 50-record database ( Figure 4B). Retrieving from a 50record database with similar query characteristics required 16 seconds while a larger 350-record database showed a 112-second retrieval time.
Discussion
In this study, we presented a proof-of-concept and use case smart contract programmed in Solidity using the Ethereum blockchain. This smart contract demonstrated the ability to store and retrieve COVID-19 patient data. The method's speed and memory per insertion were measured with varying database record sizes, demonstrating that the technique is both rapid and readily feasible. This demonstrated a practical method while also allowing for efficient memory and time usage when storing and retrieving observations. The results show that blockchain technology can be used to solve problems not specific to cryptocurrency. The results also exemplify the practicality and efficiency of smart contracts involving health-related data.
Secure maintenance of COVID-19 patient medical data is critical, as data corruption can lead to fatal medical errors and misleading public health statistics. Moreover, data should only be accessible to authorized persons such as healthcare workers and researchers. Solutions utilizing blockchain technology, such as smart contracts, could prevent centralized point-of-failure scenarios, create an immutable database, and allow decentralization, ultimately mitigating the probability of data corruption. Indeed, blockchain technology is being conceptualized in other fields as well [11][12][17][18][19][20].
During the COVID-19 pandemic, patients have presented with a wide array of symptoms. As part of diagnosing the condition, distinct criteria needed to be defined and met. Lung and infectious disease experts have noted that an especially unique factor of COVID-19 patients is the sound of their breathing, speaking, and dry coughing [21][22][23][24]. To diagnose patients' conditions and create tools that can diagnose them in physician-lacking settings, scientists have created machine-learning algorithms. These tools require large datasets as input, and their final product algorithms can rapidly and automatically function in clinics. In this case, algorithms were created to diagnose COVID-19 based on these sounds. Similarly, associations have been found between other symptoms and COVID-19. The creation of such tools is highly dependent on large datasets and their secure handling. In pandemics and other situations with large datasets containing private patient information, the Ethereum blockchain can be used to speed up secure data storage and retrieval, facilitating the faster and safer development of clinical tools.
As with any novel development, the Ethereum platform consists of limitations associated with smart contracts. For instance, programming and deploying smart contracts require knowledge of the Solidity language, which comprises numerous idiosyncrasies. Legal ramifications also need to be considered when dealing with smart contracts, since they reduce dependence on intermediaries, including lawyers. Therefore, it is imperative that all parties working with smart contracts understand the legal nomenclature of smart contract law. There is also the possibility of delayed transactions in instances where blockchain congestion results in delayed transactions and increase transaction costs. Such limitations exemplify the demand for a stabler platform. Future studies should focus on how such smart contract methods can also be applied to other problems in the medical research community. Additionally, future studies should also focus on other smart contract architectures best suited for other types of medical data. However, the potential of using smart contracts for COVID-19 information remains.
Conclusions
Blockchain technology can be used to store and retrieve COVID-19-related information securely, immutably, and rapidly. The Ethereum blockchain presents several advantages and has not yet been applied to healthcare information storage, especially during pandemics. In this study, we presented a method for securely and rapidly storing and retrieving COVID-19-related health information using smart contracts with the Ethereum blockchain. It enables the protection of patient information while facilitating more efficient research and communication.
Additional Information Disclosures
Human subjects: All authors have confirmed that this study did not involve human participants or tissue. Animal subjects: All authors have confirmed that this study did not involve animal subjects or tissue.
Conflicts of interest:
In compliance with the ICMJE uniform disclosure form, all authors declare the following: Payment/services info: All authors have declared that no financial support was received from any organization for the submitted work. Financial relationships: All authors have declared that they have no financial relationships at present or within the previous three years with any organizations that might have an interest in the submitted work. Other relationships: All authors have declared that there are no other relationships or activities that could appear to have influenced the submitted work. | 3,322.2 | 2022-01-01T00:00:00.000 | [
"Computer Science",
"Medicine"
] |
Improving Few-Shot Domain Transfer for Named Entity Disambiguation with Pattern Exploitation
,
Introduction
In order to understand a piece of text, it is often valuable to understand entities which are referred to by that text. While named entity recognition (NER) is an important aspect of this problem, a large variety of applications (e.g. financial credit risk monitoring and open source intelligence gathering) need to further connect these entities to known entities in a knowledge base (KB). This task of linking textual mentions of entities to their KB entries is known as entity linking.
Entity linking is typically further decomposed into two subtasks: candidate generation and named entity disambiguation (NED). The former is responsible for discovering a set of possible mentions of entities in a given document (for example, producing the candidates for Cambridge, MA, USA and Cambridge, UK from the surface form "Cambridge"). An NED system then takes these lists of candidates and selects which, if any, is the correct referent.
Named entity disambiguation systems achieve this by utilizing a number of pieces of information, such as related entities, the type of each candidate entity (person, location, etc.), and semantic descriptions of each entity (e.g. a snippet from the candidate's Wikipedia page). As this association is usually statistically learned from some training dataset, the performance of an NED system on a given document depends on how closely that document's domain is to that of the training data, in terms of vocabulary, syntax, and the types of entities. Because these systems are often specialized in specific domains, it is therefore necessary to curate sufficient amounts of training data for each of these applications, which is often costly.
Our work seeks to reduce the amount of data required via leveraging pretrained language models (LMs). LMs are often well-suited to assisting lowresource task setups (Tam et al., 2021), for modern language models are sufficiently powerful that their predictive distributions can be interpreted as a basic form of "common-sense reasoning".
Our contributions are as follows: (1) we take a state-of-the-art baseline NED system (Yang et al., 2019) and augment it with an additional signal from an LM fine-tuned with the ADAPET (Tam et al., 2021) procedure, which adapts language models to few-shot learning natural language processing problems. We show that this augmented system, called DCA-Prompt, achieves similar performance to the baseline in both the same and closely-related domains, but demonstrably outperforms when adapting to a new, dissimilar domain. (2) Additionally, we are releasing a new named entity linking dataset, called NEDMed, which is based on mental health news data.
Related Work
Named entity disambiguation has a storied history, stemming from the work in Bunescu and Paşca (2006), which utilized a support vector machine (SVM) (Cortes and Vapnik, 1995) kernel based on a similarity measure between input documents and the Wikipedia articles for each candidate.
Cucerzan (2007) and Kulkarni et al. (2009) were seminal works in incorporating the KB topology into this decision-making process (e.g. looking at links between candidate entities across the document), but the computational cost of these techniques was a major limitation. Numerous additional authors later provided their own approximations for this problem; a recent success in this area is known as dynamic context augmentation, or DCA (Yang et al., 2019). This technique opts to sequentially process the mentions in the document, using context related to previous extractions (the linked entity itself along with entities related to that entity) to inform subsequent extractions.
While pretrained language models have a long history (Devlin et al., 2019) and are traditionally fine-tuned using masked-language modeling in order to improve their modeling ability on new domains, pattern exploitation training (PET) (Schick and Schütze, 2021) and a densely supervised approach to pattern exploitation training (ADAPET) (Tam et al., 2021) are relatively recent applications of this fine-tuning approach. These techniques utilize the linguistic information contained in language models to solve natural language processing tasks by formulating them as cloze-style phrases. Tam et al. (2021) demonstrate its efficacy on a range of SuperGLUE (Wang et al., 2019) tasks, showing state-of-the-art or competitive performance on few-shot textual entailment and the BoolQ (Clark et al., 2019) question-answering dataset. To our knowledge, we are the first to apply pattern exploitation to NED.
The gold standard dataset for NED was defined in Hoffart et al. (2011b). This work extends the CoNLL 2003 shared task's NER dataset (Tjong Kim Sang and De Meulder, 2003) to contain links to entities from the YAGO KB (Hoffart et al., 2011a). This dataset is known in literature as the AIDA CoNLL-YAGO dataset, and is discussed further in Section 4. Additional datasets include AQUAINT (Milne and Witten, 2008), MSNBC (Cucerzan, 2007), ACE2004 (Ratinov et al., 2011), andCWEB (Guo andBarbosa, 2014). are special tokens, and items in angle brackets (<>) are substituted with information from the problem input. LCTX and RCTX are the text to the left and right of the entity mention (respectively), MENTIONTXT is the surface form of the mention, and ENTDESC is the description of the candidate under consideration. Colors represent different logical segments of each piece of the input, and [MASK] is what we prompt the language model to substitute into the input.
Methodology
In NED (Figure 1), we presume that, for a given input text containing a set of n mentions (collectively denoted M), which are the textual surface forms of entities, and a set of candidates c i (collectively, C) for each mention m i . An NED system ranks these candidates to select the one most likely to be the referent entity. A full formal description of the task is given in Appendix A.
ADAPET solves natural language understanding tasks by "filling-in-the-blank" in natural language patterns. To solve a downstream task, such as classification or question-answering, one formulates the problem instance as some form of prose containing a masked token. A fine-tuned language model is then used to infer which word makes the most sense to substitute this masked token. In order to apply ADAPET to NED, it was necessary to formulate an appropriate pattern. Our system fine-tunes the HuggingFace (Wolf et al., 2020) bert-large-uncased model using the pattern shown in Figure 2, treating the task as a binary classification problem (answering, "does this candidate link to this mention?"). The develop- Figure 3: Few-shot learning curves for ADAPET-NED, DCA, and DCA-Prompt models. "DCA (local only)" measures performance on the DCA model without entity context enabled, which is more directly comparable to the ADAPET-only line. The x-axis indicates the number of mentions of aida-train used while training. ment process through which we chose this base model and pattern is described in Appendix B. As done in the codebase provided by Tam et al. (2021), we fine-tune the model to produce the word "true" for correct mention-candidate pairs and the word "false" for incorrect ones. The entity description ([ENTDESC]) may be something like a paragraph from an encyclopedia entry or a verbalized form of its knowledge base relationships, as described in Mulang' et al. (2020). This representation should capture enough context to uniquely identify the entity. In our work, we use the first paragraph of the candidate's Wikipedia page.
In order to accurately judge this approach versus state-of-the-art NED systems, we additionally augment an existing state-of-the-art baseline system with an input signal from our ADAPET-tuned LM. The baseline system we selected is based on DCA, and is described by Yang et al. (2019). In short, this model, is a feedforward neural network which accepts features based on the similarity of the mention context and trained entity embeddings (Ψ C ), entity type information (Ψ T ), and coherence between the candidates and both previously linked entities and entities linked to those entities (Φ and Φ ′ ). For further information, readers are referred to Yang et al. (2019). Further technical details regarding our training setup can be found in Appendix C.
Our research focused on the following questions: RQ1. Does using ADAPET for NED perform similarly to the baseline cross-domain scenarios for similar domains?
RQ2. Does using ADAPET for NED reduce the amount of data required to learn the task?
RQ3. Does using ADAPET for NED improve the ability for a NED system to transfer to another dissimilar domain?
Results
To answer our first two research questions, we compared ADAPET-NED (ADAPET trained on our NED prompt) against three versions of the baseline DCA system: one which operates as normal (ETZH-Attn+DCA-SL from Yang et al. (2019)), denoted "DCA", one which only has "local" features enabled (i.e. no coherence-based features, which makes it more directly comparable to ADAPET-NED), denoted "DCA (local only)", and a version of the DCA system which has been augmented to receive the output from the ADAPET model's "yes" prediction as an input, denoted "DCA-Prompt". We model the (2-way k-shot) few-shot learning curves of these models in Figures 3(a) and 3(b), evaluating on in-domain datasets (training on different sized subsets of AIDA's aida-train training split and evaluating on aida-A and aida-B development and evaluation splits). As expected, this graph shows that ADAPET-NED greatly underperforms the baseline system. To explore RQ1, we measured cross-domain performance on various publicly available datasets and compare to existing benchmarks. Specifically, we looked at F1 performance on MSNBC, AQUAINT, ACE2004, and CWEB (described in Section 2), as done in Yang et al. (2019). These datasets are all general knowledge corpora, based on either news or encyclopedia pages. We train our models on aidatrain and evaluate across all datasets. The results are shown in Table 1. We see that, while ADAPET-NED is not competitive, augmenting DCA with features from ADAPET (DCA-Prompt) yields performance ranging from comparable to superior, with state-of-the-art performance on the ACE2004 dataset. Additionally, for RQ2, we find that all three of these models have similar data requirements. We quantify this by using the Kneedle al- gorithm (Satopaa et al., 2011) to locate the "knees" in each of the curves in Figure 3, which showed diminishing returns at around 1,000 mentions for all three.
In order to assess RQ3 in a real-world context, we adapted our trained models to a dataset tailored to the medical domain, which is quite different from AIDA's general news domain. We created a dataset, denoted NEDMed 1 , containing 110 internet articles on mental health news, which were partitioned into 66 training documents (NEDMedtrain) and 44 evaluation documents (NEDMeddev), containing 2,839 and 1,841 mentions, respectively. Documents were manually annotated for person, location, and organization types, along with a variety of others. For a full list of types and further details on this dataset, see Appendix D. For our experiments, we only utilize entities which have Wikipedia links (4,342 mentions, or roughly 92% of the total 4,680). Table 2 and Figure 3(c) describe the results of the baseline DCA system, our ADAPET-NED and DCA-Prompt systems on this data. The NEDMeddev scores on models trained with aida-train (the first group in Table 2) represent zero-shot (crossdomain) scenarios. The models trained on the combined data represent transfer learning scenarios in which we tune an AIDA-trained model on NEDMed data. We additionally report scores on aida-B in order to monitor catastrophic forgetting. Finally, to measure the contribution of aida-train, we trained models using NEDMed-train alone, and found lower NEDMed-dev scores across the board (a roughly 3% drop in F1).
We find that our DCA-Prompt system yields superior performance in both the zero-shot and transfer learning scenarios. Notably, the zero-shot performance of DCA-Prompt is higher than all three 1 The NEDMed dataset is available to download at https: //github.com/basis-technology-corp/ NEDMed . metrics of the baseline DCA system.
Conclusions and Future Work
These results indicate that pattern exploitation training can effectively be utilized for named entity disambiguation. While results are not state-of-the-art when used in isolation, combining an ADAPETbased classifier with an existing model which can incorporate global context, such as DCA, improves the capacity of that model to flexibly adapt to data from different domains. Our new NEDMed dataset both provides evidence for this and represents a new domain-specific benchmark which can be used by future NED research.
There are a number of ways in which this work could be built upon in the future. This work focused on shifts in domain related to the documents in which mentions are extracted from. Another important type of domain shift relates to large changes in the underlying KB. While we expect the system would be able to adapt, this has not been quantified. Additionally, the optimal strategy for designing patterns for use with ADAPET-style techniques is still an open research question (Liu et al., 2021); as this work relied on human-produced patterns, it is certainly possible that accuracy or data requirements could be improved with more clever pattern design.
Risks and Limitations
The authors of this paper believe that this work does not introduce any unique risks or limitations; however, we shall note some which are inherent to named entity disambiguation in general. As it is a central inspiration of this work, one of the most noteworthy limitations is that of cross-domain applicability. That is, the performance of our NED system on a given datum remains a function of how closely that datum reflects the data upon which the system was trained. While our work narrows the gap in performance, it remains the case that data Table 2: Transfer learning performance on aida-B and NEDMed-dev datasets when trained on aida-train ("AIDA") and NEDMed-train. Best scores and our systems are in bold. "ADAPET" denotes our ADAPET-NED system.
from extremely different domains (e.g. a different KB which is dissimilar from Wikipedia) will not be linked as accurately as data from the same domain.
The primary societal risk of NED systems is that of surveillance. While it does not increase the ability to collect data which may pertain to a given entity, well-performing NED systems reduce the amount of human labor which is needed to filter through false positives returned by data collection streams. This reduces the total amount of effort required for organizations to precisely aggregate information about specific entities across large quantities of data.
A Formal Description of Named Entity Disambiguation
In this section, we provide a formal description of the terminology of named entity disambiguation, building upon the brief outline in Section 1. These terms are shown in Figure 1. Recalling from Section 3, in NED, we presume that, for a given input text, a candidate generator Table 3. "(yes/no)" and "(true/false)" indicate the values used for the "[MASK]" token. The x-axis indicates the number of mention-candidate pairs used during training.
detects n mentions in the document (denoted M), which are the textual surface forms of entities in the document which may or may not link to a KB entity (much like the entities extracted by a named entity recognition system). A list of candidates c i (collectively denoted C) is associated with each mention m i in the document. The goal of a named entity disambiguation system is to inspect the output from the candidate generator and determine the correct candidate for each mention. Typically, this is done via scoring each mention-candidate pair. Mathematically, this is done with a scoring function s as follows: Note that s is able to incorporate arbitrary amounts of context in its decision-making process (e.g. other mentions, candidates of other mentions, etc.). In both most recent and this work, s is a neural network trained via gradient descent.
B Pattern and Transformer Analysis
Before we could answer our research questions, it was necessary to understand which choice of pattern makes the most sense for this task. To this end, we needed to first produce a set of patterns to choose from. Unlike many of the SuperGLUE tasks, there was not an obvious choice for what a good pattern may be, so we experimented with a few, shown in Table 3; nonetheless, this served as a good exercise in how to design patterns for use with these systems, which should prove helpful for others. Note that we did experiment with prompt tuning (Lester et al., 2021), but this did not give good results. First, there was a choice of whether to create a binary or n-ary pattern (i.e. "is this the correct candidate?" vs. "which of these are the correct candidate?"). Our work uses the former, as some preliminary empirical results from the latter yielded poor results. For such binary patterns, we need to include two pieces of information: a snippet from the input document (the mention, along with its surrounding context), and some sort of information about the candidate that we want to evaluate. Additionally, these two pieces of information need to be bridged together by the pattern in such a way that we have a masked token which can be filled in to answer the "is this the correct candidate?" question. One notable aspect of these patterns was the decision to utilize a distinct "[MENTION]" token inside of the first pattern, in place of the surface form of the mentioned entity. This is done in order to more directly relate the in-context mention to the question at the end of the pattern, as the candidate description presumably contains many instances of the mention's surface form. To represent the candidates in the model inputs, we rely on the existence of textual descriptions of each entity (ENTDESC).
To determine which would be optimal, we train each of the three patterns from Table 3 on aida-train in order to compare and contrast two pieces of information: the overall accuracy on aida-A when using each pattern and which of these patterns required the least amount of training data to achieve this accuracy. As done in Tam et al. (2021), we tune the transformer for a single epoch over the data, and we sample the aida-A performance every 640 candidate-mention pairs. As is standard, we evaluate aida-A using in-KB accuracy, which is simply the accuracy for all aida-A mentions whose correct answer is in the knowledge base. The results of this analysis are shown in Figure 4. As mentioned in Section 3, we fine-tune various HuggingFace Transformers (Wolf et al., 2020) models to produce the word "yes" or "true" for correct mentioncandidate pairs and the word "no" or "false" for incorrect ones. We find that all patterns perform roughly the same (whether using "yes" and "no" or "true" and "false" as the pattern output), with the exception of P3 with a "true"/"false" output (which performs slightly worse). We additionally experimented with ensembling the three patterns together, but this yielded performance worse than using patterns in isolation. As it yielded the greatest overall Figure 5 (smoothing was done by averaging each data point with its immediate neighbors). Average in-KB accuracy is the mean accuracy across all points in training (higher values indicate better few-shot learning ability). score, for the analysis in Section 4, we utilize P2 with "true"/"false" as our ADAPET pattern. Furthermore, we needed to understand which pretrained language model would provide the best performance on this task when fine-tuned with the ADAPET training procedure. To this end, we trained a number of models with different pretrained transformers, with the results in Figure 5 and Table 4. While they all largely converged to a similar in-KB accuracy on the aida-A dataset, the bert-large-uncased model reached this value more rapidly than the other transformer models (quantified by its average accuracy), so it was ultimately chosen for the experiments in Section 4. Notably, all of these models other than Longformer (Beltagy et al., 2020) accept inputs up to a maximum length, so our inputs were trimmed to a maximum length of 256 (the trimming was done in a manner balanced across the color-coded segments of Table 3, with the segment contain the "[MASK]" token not being truncated; this strategy is roughly equivalent to that which is used in Tam et al. (2021)). We note that this is well above the average length of inputs. Additionally, we investigated Longformer as an means of reducing the amount of truncation required, but increasing the length did not yield any noticeable improvement in overall performance over bert-large-uncased.
C Training Details
For the ADAPET model, as described in Appendix B we use HuggingFace's bert-large-uncased (Lan et al., 2020) model as our base model. Each ADAPET input was truncated to a length of 256, and a batch size of 16 is used for the gradient updates. We use the AdamW optimizer (Loshchilov and Hutter, 2019) with a learning rate γ = 10 −5 and weight decay of λ = 10 −2 . The learning rate is updated according to a linear scheduler.
For the DCA model, we adopt the same hyperparameters as Yang et al. (2019), using the Adam (Kingma and Ba, 2015) optimizer with a learning rate of γ = 2 × 10 −4 . To limit the scope of these experiments, we focus on the best-performing DCA configuration, which is based on a supervised learning training strategy (referred to in the original paper as DCA-SL), with mentions ordered by offset. Reported scores are the best performance measured over up to 2 500 epochs.
For the combined DCA-Prompt model, we first train the ADAPET model on the dataset and then feed its outputs into the DCA model during a separate training session. This effectively means that we train the ADAPET model and freeze its weights when training the final DCA-Prompt model. Future work will aim to model these two components end-to-end.
Experiments were run using a single NVIDIA Tesla T4 GPU on a Google Cloud Platform n1-standard-8 machine.
The ADAPET model takes roughly 18 hours to fully train and evaluate for a single pattern. The DCA model takes roughly four hours to fully train and evaluate.
C.1 Dataset Information
For the bulk of our baseline experiments, we utilize the AIDA CoNLL-YAGO NED dataset (Hoffart et al., 2011b), as provided by Yang et al. (2019). This dataset is split into three pieces: aida-train, containing 18,448 mentions across 942 documents; aida-A, containing 4,791 mentions across 216 documents and typically used as a development set; and aida-B, containing 4,485 mentions across 230 documents and typically used as an evaluation set. Each item from this version of the dataset consists of a mention and a list of candidate Wikipedia entities. Less than 1% of the mentions in aida-A and aida-B do not include the correct candidate in their lists; as with Yang et al. (2019)'s work, these are skipped when evaluating models.
D NEDMed Datasheet
This datasheet template is taken from Gebru et al. (2021).
Motivation
For what purpose was the dataset created? Was there a specific task in mind? Was there a specific gap that needed to be filled? Please provide a description. The goal was to create an English named entity linking dataset based on health-related text.
Who created this dataset (e.g., which team, research group) and on behalf of which entity (e.g., company, institution, organization)? This annotated dataset was produced by BasisTech.
Who funded the creation of the dataset? If there is an associated grant, please provide the name of the grantor and the grant name and number. The production of this dataset was funded by Ba-sisTech.
Any other comments?
Composition What do the instances that comprise the dataset represent (e.g., documents, photos, people, countries)? Are there multiple types of instances (e.g., movies, users, and ratings; people and interactions between them; nodes and edges)? Please provide a description. The items in the dataset are documents annotated with metadata.
How many instances are there in total (of each type, if appropriate)? There are 110 documents in the dataset. Of which, 66 comprise NEDMed-train and 44 comprise NEDMed-dev. The following is the breakdown of mention types in NEDMed-train: And the following is the breakdown for NEDMed-dev: Does the dataset contain all possible instances or is it a sample (not necessarily random) of instances from a larger set? If the dataset is a sample, then what is the larger set? Is the sample representative of the larger set (e.g., geographic coverage)? If so, please describe how this representativeness was validated/verified. If it is not representative of the larger set, please describe why not (e.g., to cover a more diverse range of instances, because instances were withheld or unavailable).
What data does each instance consist of? "Raw" data (e.g., unprocessed text or images) or features? In either case, please provide a description.
Each instance is a text document that includes metadata with information such as source, publication date, and language. Entities in the document have been annotated with character offsets, knowledge base identifiers, and types. The possible types according to annotation guidelines were Location, Organization, Person, Product, Nationality, Religion, Title, Disease, Symptom, Substance, and Treatment. A (visually rendered) example of an annotated sub-section of a document is the following (brackets have been placed around annotated entities, and entities with the same color represent ones which are linked to the same Wikipedia entity): Is there a label or target associated with each instance? If so, please provide a description. Each document in the dataset contains entity mentions, which are associated with a knowledge base identifier (either a Wikidata QID or a custom knowledge base ID) and an entity type.
Is any information missing from individual instances? If so, please provide a description, explaining why this information is missing (e.g., because it was unavailable). This does not include intentionally removed information, but might include, e.g., redacted text.
No.
Are relationships between individual instances made explicit (e.g., users' movie ratings, social network links)? If so, please describe how these relationships are made explicit.
N/A
Are there recommended data splits (e.g., training, development/validation, testing)? If so, please provide a description of these splits, explaining the rationale behind them. Yes. The dataset is split into a training dataset (NEDMed-train) and a development/evaluation dataset (NEDMed-dev). This was done by randomly partitioning the documents.
Are there any errors, sources of noise, or redundancies in the dataset? If so, please provide a description. There is possible noise in the dataset. All data was annotated manually, but the final Krippendorff's α value (pairwise inter-annotator agreement) for the NER annotations was 0.768 and for the linking annotations was 0.767. This means that there remained some level of disagreement among the annotators, which could manifest as noise in the data.
Is the dataset self-contained, or does it link to or otherwise rely on external resources (e.g., websites, tweets, other datasets)? If it links to or relies on external resources, a) are there guarantees that they will exist, and remain constant, over time; b) are there official archival versions of the complete dataset (i.e., including the external resources as they existed at the time the dataset was created); c) are there any restrictions (e.g., licenses, fees) associated with any of the external resources that might apply to a future user? Please provide descriptions of all external resources and any restrictions associated with them, as well as links or other access points, as appropriate.
The linking information in the dataset refers to Wikidata entities. These are publicly available without restriction and will not change.
Does the dataset contain data that might be considered confidential (e.g., data that is protected by legal privilege or by doctor-patient confidentiality, data that includes the content of individuals non-public communications)? If so, please provide a description.
No.
Does the dataset contain data that, if viewed directly, might be offensive, insulting, threatening, or might otherwise cause anxiety? If so, please describe why.
Yes. Some articles mention sensitive mental health topics such as suicide.
Does the dataset relate to people? If not, you may skip the remaining questions in this section.
Not directly. The articles in the dataset relate to mental health; these may be news stories involving people.
Does the dataset identify any subpopulations (e.g., by age, gender)? If so, please describe how these subpopulations are identified and provide a description of their respective distributions within the dataset.
N/A
Is it possible to identify individuals (i.e., one or more natural persons), either directly or indirectly (i.e., in combination with other data) from the dataset? If so, please describe how.
N/A
Does the dataset contain data that might be considered sensitive in any way (e.g., data that reveals racial or ethnic origins, sexual orientations, religious beliefs, political opinions or union memberships, or locations; financial or health data; biometric or genetic data; forms of government identification, such as social security numbers; criminal history)? If so, please provide a description. N/A Any other comments?
Collection Process
How was the data associated with each instance acquired? Was the data directly observable (e.g., raw text, movie ratings), reported by subjects (e.g., survey responses), or indirectly inferred/derived from other data (e.g., part-of-speech tags, modelbased guesses for age or language)? If data was reported by subjects or indirectly inferred/derived from other data, was the data validated/verified? If so, please describe how.
The document text was collected by searching the sites https://theconversation.com/ and https://en.wikinews.org/ for articles related to mental health. Each document was then annotated by a minimum of two human annotators. In the cases where the annotators disagreed, an adjudication process was used to determine the final set of annotators.
What mechanisms or procedures were used to collect the data (e.g., hardware apparatus or sensor, manual human curation, software program, software API)? How were these mechanisms or procedures validated? The unannotated data was collected manually by two employees of the BasisTech data team. Annotators were provided with a set of instructions describing how to annotate for named entities and their links. Annotation was done using an internal proprietary NLP annotation tool, which allows metrics such as inter-annotator agreement to be measured across an annotation project.
If the dataset is a sample from a larger set, what was the sampling strategy (e.g., deterministic, probabilistic with specific sampling probabilities)? The documents were chosen by hand based on their content and metadata in order to target news topics related to health and mental health.
Who was involved in the data collection process (e.g., students, crowdworkers, contractors) and how were they compensated (e.g., how much were crowdworkers paid)?
The collection of the raw data was performed by employees of BasisTech. Annotation was performed by three experienced Israeli contractors with whom Basis had worked with prior and were compensated at $15-30 per hour. One contractor was a native English speaker, and the other two were native Hebrew speakers with high levels of English competency. The arbitration process for annotation conflicts was performed by a BasisTech employee who is a native English speaker.
Over what timeframe was the data collected? Does this timeframe match the creation timeframe of the data associated with the instances (e.g., recent crawl of old news articles)? If not, please describe the timeframe in which the data associated with the instances was created.
The unannotated documents were collected between August 5th, 2020 through August 10th, 2020. The original publication of the documents ranged from April 5th, 2005 through May 29th, 2020.
Were any ethical review processes conducted (e.g., by an institutional review board)? If so, please provide a description of these review processes, including the outcomes, as well as a link or other access point to any supporting documentation.
No.
Does the dataset relate to people? If not, you may skip the remaining questions in this section.
No.
Did you collect the data from the individuals in question directly, or obtain it via third parties or other sources (e.g., websites)? N/A Were the individuals in question notified about the data collection? If so, please describe (or show with screenshots or other information) how notice was provided, and provide a link or other access point to, or otherwise reproduce, the exact language of the notification itself.
N/A
Did the individuals in question consent to the collection and use of their data? If so, please describe (or show with screenshots or other information) how consent was requested and provided, and provide a link or other access point to, or otherwise reproduce, the exact language to which the individuals consented.
N/A
If consent was obtained, were the consenting individuals provided with a mechanism to revoke their consent in the future or for certain uses? If so, please provide a description, as well as a link or other access point to the mechanism (if appropriate).
N/A
Has an analysis of the potential impact of the dataset and its use on data subjects (e.g., a data protection impact analysis) been conducted? If so, please provide a description of this analysis, including the outcomes, as well as a link or other access point to any supporting documentation. N/A Any other comments?
Preprocessing/cleaning/labeling Was any preprocessing/cleaning/labeling of the data done (e.g., discretization or bucketing, tokenization, part-of-speech tagging, SIFT feature extraction, removal of instances, processing of missing values)? If so, please provide a description. If not, you may skip the remainder of the questions in this section. The documents were tokenized with BasisTech's Rosette® Text Analytics linguistic analysis software before annotation (entity mention annotations align with token boundaries).
Was the "raw" data saved in addition to the preprocessed/cleaned/labeled data (e.g., to support unanticipated future uses)? If so, please provide a link or other access point to the "raw" data. The original text of each collected document is included in each instance in the dataset.
Is the software used to preprocess/clean/label the instances available? If so, please provide a link or other access point. No, proprietary software was used.
Any other comments?
Uses
Has the dataset been used for any tasks already? If so, please provide a description.
Yes, this paper.
Is there a repository that links to any or all papers or systems that use the dataset? If so, please provide a link or other access point.
Not at present.
What (other) tasks could the dataset be used for? Named entity recognition (NER).
Is there anything about the composition of the dataset or the way it was collected and preprocessed/cleaned/labeled that might impact future uses? For example, is there anything that a future user might need to know to avoid uses that could result in unfair treatment of individuals or groups (e.g., stereotyping, quality of service issues) or other undesirable harms (e.g., financial harms, legal risks) If so, please provide a description. Is there anything a future user could do to mitigate these undesirable harms?
Not to our knowledge.
Are there tasks for which the dataset should not be used? If so, please provide a description.
No.
Any other comments?
Distribution
Will the dataset be distributed to third parties outside of the entity (e.g., company, institution, organization) on behalf of which the dataset was created? If so, please provide a description.
Yes. The data shall be publicly released alongside this paper.
How will the dataset will be distributed (e.g., tarball on website, API, GitHub) Does the dataset have a digital object identifier (DOI)? | 8,086.4 | 2022-01-01T00:00:00.000 | [
"Computer Science"
] |
Cross Correlation for Condition Monitoring of Variable Load and Speed Gearboxes
The ability to identify incipient faults at an early stage in the operation of machinery has been demonstrated to provide substantial value to industry. These benefits for automated, in situ, and online monitoring of machinery, structures, and systems subject to varying operating conditions are difficult to achieve at present when they are run in operationally constrained environments that demand uninterrupted operation in this mode. This work focuses on developing a simple algorithm for this problem class; novelty detection is deployed on feature vectors generated from the cross correlation of vibration signals from sensorsmounted on disparate locations in a power train. The behavior of these signals in a gearbox subject to varying load and speed is expected to remain in a commensurate state until a change in some physical aspect of the mechanical components, presumed to be indicative of gearbox failure. Cross correlation will be demonstrated to generate excellent classification results for a gearbox subject to independently changing load and speed. It eliminates the need to analyze the highly complex dynamics of this system; it generalizes well across untaught ranges of load and speed; it eliminates the need to identify and measure all predominant time-varying parameters; it is simple and computationally inexpensive.
Introduction
The dynamics of the vibrations generated by a gearbox subject to changing load and speed are complex and nonlinear.Faults in bearings, gears, or other aspects of prime movers can easily be masked by the effects of these state changes alone when one fails to consider their effects on decision rules.The detection of faults in this class of machineries is a growing concern in the literature.In this work, we adapt a technique from sensor failure analysis to reduce this present problem's complexity.A common approach in detecting failure in sensors employs decision rules based on the cross correlation of their signals; in broaching this technique to variable-state machinery, the authors note that vibrations at disparate locations in a power train should be correlated to one another (e.g., the spectra of vibrations from the output shaft of a gearbox are related to those of the input shaft by the gear ratio of the gearbox).Signals from disparate locations of a power train may contain similar vibration from components along the train; for instance, the load on the gearbox's bearings is modulated by the meshing of the gear's teeth and its vibrations or acoustics will be apparent at both the input and the output of the gearbox (and possibly at more distant locations in the train; see [1]).The cross correlation signal between these vibration signals should remain commensurate until components of the train change-a state presumed indicative of faults.
Under this hypothesis, the authors propose deploying standard novelty detection on feature vectors generated from the cross correlation signal generated between disparate vibration sensors.Past efforts by the authors focused on adapting either novelty-detection techniques or feature vectors in order to address this problem.These algorithms required the investigators to measure all predominant state parameters and to include them in the algorithm [1,2].While the proposed techniques were shown to work well, they suffered from various limitations.Some classification schemes work only for one changing system input parameter [1].Others require measurement of a gearbox's load which can be either a costly or cumbersome requirement when an inline load cell needs to be installed on a system not fitted with it.Finally, the computational complexity of others requires large processing facilities not typically available on distributed embedded systems employed in condition monitoring.The cross correlation technique should eliminate or mitigate all of these drawbacks.This approach should provide an excellent means of failure detection in systems whose dynamics are too complex for traditional approaches and consequently may extend well beyond the monitoring of variable load and speed gearboxes.
To validate these conclusions, the necessary theoretical background is first explored including a review of cross correlation and how it is presently employed in this field as well as an overview of other existing approaches for solving this class of problems.The underlying methodology is subsequently described, from a description of the employed mechanical test bench to the details of each of the steps in the classification problem.Finally, the results are demonstrated to establish the flexibility of this simple approach.
Background
The mathematics of cross correlation is first reviewed followed by an overview of related existing techniques.
Cross Correlation Analysis.
Cross correlation analysis provides a signal representing the measure of the similarity between two signals as a function of time lag , defined as where ⊗ denotes the cross correlation function and * denotes complex conjugation; similarly, it can be expressed in discrete form It is used extensively in pattern-recognition for speech, fingerprint and face recognition, automatic target recognition, and so forth.In these applications, typically one cross correlates a reference pattern with a test pattern when the two patterns are expected to lack shift invariance.The cross correlation signal between two patterns will have a peak at the shifted value if they have some similarity.
Cross Correlation of Systems Subject to Common Excitation.
In this work, signals from disparate aspects of machinery, under common excitation, are cross correlated in order to simplify discerning the system's health when the excitation is nonstationary.If two linear systems, with impulse response functions 1 () and 2 (), are commonly forced with some function (), having equivalent frequency domain representation of ( 0 ), the particular solutions for the systems' response will be the product of the forcing function and the system's impulse response for all in ( 0 ); that is, 1 = 1 ( 0 ) * ( 0 ) and 2 = 2 ( 0 ) * ( 0 ).From elementary Laplace and Fourier transform theory, it is known that the frequency domain representation of the convolution of two signals is the product of their frequency domain representations.Cross correlation is equivalent to the convolution operation except without the folding operation; as such the frequency domain representation o the cross correlation of two signals is the product of the two signal's frequency domain representations.Since the linear systems are forced with the same function, their output signals' bandwidths overlap and the frequency domain representation of the cross correlation operation returns a product of the two system's impulse response functions and the forcing function.The impulse response function for each system is determined by the system's parameters (e.g., for a spring, the impulse response function is a function of the spring's stiffness, the damping constant, etc.).The cross correlation of the two systems' output therefore is a relation given by the system's parameters.If any parameters of a system change, the cross correlation of the two systems' output will change; it is on this basis that this work is advanced.
The vibration from a gearbox is inherently nonlinear and some of the assumptions of the foregoing therefore break down.Complex pattern-recognition techniques like novelty detection are engaged to handle these aspects.
Sampled systems are discrete in nature which was not presumed in the above analysis.The discrete systems under scrutiny herein are made discrete by sampling the continuous phenomena.The argumentation of the above is very similar in discrete form and a direct analogy can be made between the transforms of discretized form and the continuous form.
Relevant Cross Correlation Techniques from the Literature.
Cross correlation is used heavily in signal processing for denoising purposes.Several examples of denoising in the domain of fault detection can be found in the literature; in [3], the authors used cross correlation from two proximate vibration sources for signal-to-noise ratio improvement while [4] used cross and autocorrelation for denoising.The authors in [5] exploited the auto-and cross correlation of different variables for signal processing in developing a fault-detection technique.
Cross correlation is used in a similar vein as the present approach in the detection of failed sensors as was the case in [6] whose authors used cross correlation between two flow sensors along with neural networks to verify sensor accuracy.The work in [7] acknowledges the dynamic nature of a motor run by an adjustable speed drive and the resultant effects on monitored signals are one of the common factors that yield erroneous fault tracking and unstable fault detection; the authors employed matched filtering (i.e., cross correlation between expected fault signals and actual motor current signals) the result of which is fed through a statistical hypothesis-testing fault-detection regime.Statistical-process monitoring with spectral clustering was used to classify samples according to differences in correlation among measured variables in [8].In [9] cross correlation of the fault-response echo in electrical-power transmission systems from testinput excitation was used to detect potentially faulted cables.
Jiang et al. [10] used the correlation dimension (a type of fractal dimension) in gearbox fault diagnosis.
More directly related techniques can be found in a number of other works.For instance, Parlar employed a similar methodology to that of this thesis in the monitoring of vibrating screens in [11].In [12] Napolitano et al. exploited cross correlation of an airplane's pitch and yaw state variables along with neural networks for fault identification in airplane systems.Rajamani et al. found the cross correlation between healthy and faulted transformer winding signals that was used to generate statistical feature vectors for classification [13].In [14], Wu and Sun used the cross correlation of energy performance of a variable-air-volume (VAV) unit in an HVAC system [15] and the outside temperature as the criteria to evaluate the VAV health.
Cross correlation is used heavily in this field but the methodology proposed herein on this particular class of problems does not appear to exist in the literature.
Established Techniques.
In the literature, there are a number of other algorithms focused on means other than correlation based fault detection for this complex class of machineries.Nonlinear principal-component analysis (NLPCA) in [16], advanced signal processing in [17,18], adaptive filters in [19,20], and adaptations to pattern-recognition techniques in [21][22][23][24] are all well established-each having differing strengths and weaknesses.
To provide a baseline for comparison for the approach advanced within, a comparison between a number of related techniques developed by the present authors will be undertaken.In [1] the authors explored expansions to the work by Worden et al. in [25]; Worden et al. suggested that vibration data from structures be grouped into discrete ranges of the time-changing parameters whose statistics (mean and covariance) are regressed or interpolated to develop a health rule as a function of the time-varying parameters.The work in [1] applied this approach to data from real gearbox vibrations along with an augmentation to Worden's approach that focused on first whitening the statistical distribution so that any variant of novelty detection could be employed.Both techniques were subject to the assumption of normally distributed data and the double curse of dimensionality, a phenomenon occurring when there is a need not only to gather sufficient data to describe a complex high-dimensional problem space but also to do so for continuous changes in that problem space (e.g., in the form of changing speed or load).These initial investigations were conducted with only one time-changing parameter; in this work, two timechanging parameters are used (i.e., speed and load).While a large amount of data has been collected (nearly 20,000 feature vectors generated with ambitious segmentation), they are insufficient to accurately characterize the behavior of the gearbox with these approaches due to the double curse of dimensionality.
In an upcoming work, the present authors suggested the almost trivial approach of adding a gearbox's average speed over a feature vector's segment to that feature vector.The results generated with the same experimental data were found to be excellent; unfortunately, the fault-detection methodology does not extend beyond one time-varying parameter.The confusion eliminated in adding one time-varying parameter to the feature vector is again reintroduced when another time-varying parameter is added.
In a different upcoming work, the authors suggest using the parameters of a discrete state-space model as elements of the feature vector in the novelty-detection problem [26].In a simple view, this state-space model can be regarded as the transfer function of a gearbox modeled as a torsional spring; the state-space model's parameters are ultimately functions of the physical nature of the gear (i.e., stiffness, damping, geometric configuration, etc.).These parameters ought to be insensitive to changes in load and speed and should be highly indicative of incipient fault states.The model is generated by assuming that the gearbox's input speed and load are the inputs to a MIMO system; the vibration signal at any point on the machine is used as the output signal and the MIMO model formed with ARMAX techniques [27].While the vibration problem being modeled with this linear state-space approach is in reality nonlinear, the use of novelty detection to develop a boundary around a set of these linear models is shown to provide adequate adaptation to the underlying nonlinear problem.The approach was shown to eliminate the double curse of dimensionality and assumption of normally distributed data.As evidence of the model's sound nature, the results demonstrated excellent generalization to speeds and loads not experienced during training.The only limitations to the approach are the need to collect speed and load signals (a potentially costly consideration) and the computationally intensive nature of the algorithms for generating these models.
Experimental Configuration
This work focuses on the use of the parameters generated by cross correlating signals from sensors on disparate components of a machine.The pattern-recognition problem as advanced by [28] focuses on first collecting and conditioning signals (in this case, on a simulation test bench), segmenting them, and transforming them into -dimensional feature vectors that are ultimately fed into pattern-recognition solutions.The steps for this problem instance are described below.
Apparatus.
The fault-detection algorithm proposed herein was evaluated based on data collected from a gearbox under realistic load and speed as shown in Figure 1.The test bench is described in further detail in [29].
The gearbox's independent load and speed profiles were affected via a 25 hp and 50 hp AC induction servomotor ultimately controlled by two Baldor variable frequency drives (VFDs) with appropriate capacity.This gearbox was a singlestage reduction spur gearbox from SpectraQuest.Its shaft was supported by Rexnord ER-10 deep-groove rolling element ball bearings.Coupling between the motors and gearboxes was achieved through a combination of rigid shaft couplings and two zero-backlash alignment-enhancing R + W BK3 Bellows flexible couplings.The entire drive train from the load to speed motors is shown in Figure 2. Control and data acquisition were achieved primarily with a national instruments (NI) PCIe-7851-R field programmable gate array (FPGA) card with 8 channels of analog input/output and 96 channels of digital input/output.The control and data acquisition routines were written in Lab-VIEW code for both the real-time Windows PC and mounted FPGA card (capable of loop iteration in the nanosecond range).This PC was further fitted with an NI PCI-4472 card supporting 8 channels of IEPE acceleration data.
Four accelerometers, sampled at 10 kHz, were fitted on diverse components of the drive train.One accelerometer was mounted radially on the bearing of the drive motor, two were mounted radially and orthogonally to one another on the output side of the gearbox near the input shaft, and the final accelerometer was mounted on the input side of the gearbox near the output shaft.A Lorenz Messtechnik DR-2112-R inline torque meter was fitted on the input side of the gearbox and data were collected from it at 1 kHz with the FPGA card.Tachometer signals from the two motors were first counted by sampling the TTL pulses at a rate of 40 MHz on the FPGA card; this count signal was then sampled at 10 kHz and written to disc.
Control is achieved by using two analog output lines, one to each of the motors VFDs.A typical speed/load profile employed during data collection is shown in Figure 3.
Faulted Components.
The first data set consisted of spur gears with a gear ratio of 3 : 1 in a reduction arrangement.Data were collected by swapping healthy and faulted components; bearing faults consisted of rolling elements with rough balls, a chopped ball, and inner and outer race faults of varying severity.Faulted gear signals were generated through the use of eccentric gears and two different gears with increasing rootcrack depth (generated by wire electric discharge machining).
An additional set of gears consisting of a ratio of 80 : 48 were deployed in order to show the effect of the analyzed techniques on a different set of interesting gear faults including a gear with both a missing tooth and crack as well as a gear with teeth with progressively less material.
Signal Segmentation.
Feature vectors are generated from continuously sampled signals split into meaningful and coherent intervals.In selecting the size of a signal segment, one must ensure that there is sufficient data to confirm that all necessary mechanical behavior is captured and that subsequent segments ensure a coherent comparison (i.e., each segment should accurately represent the cycle of mechanical behavior).When monitoring systems experience changes in state, the problem can become slightly more complex.One must gather sufficient data to adequately characterize the feature in question; there might also be a need to minimize the duration of the interval in order to eliminate large changes in signal behavior due to changes in system states.This is particularly true where the feature vectors are sensitive to the changing states and other means of ensuring accurate classification are employed (see [1]).
The constraints on segmentation in the problem at hand are more similar to the steady-state system case.Since the objective is to seek parameters immune to changes in system state with cross correlation based feature vectors, the only concern is the coherence and sufficiency of the segment.Consequently, concerns over accelerations and higher level rates of change in a segment from state variables, such as speed, should provide little impact.These constraints will be satisfied by using a variable-length period with a fixed number of shaft rotations (i.e., 15).
Feature Vectors.
Feature parameters are formed from processing signal segments and are combined together to form an -dimensional vector.The authors' favored approach in the form autoregressive (AR) models will be considered; AR models provide a high-dimensional feature vector by minimizing a signal in the least-squares sense to the most representative samples (the parameters of these models have a strong tie to the frequency characteristics of the signal) (see [30] for a better background).
Pattern-Recognition Algorithm.
In developing a model of a machine's behavior, it is generally only a simple task to collect data representative of the machine's healthy state; collection of data from faulted states is either too difficult because of the varied number of such states or economically/operationally infeasible to do so (particularly with machinery in use in industry).This class-imbalance problem is typically resolved through the use of novelty detection where a decision boundary is fit around exemplars of dimensional vectors derived from system signals that ideally represent the healthy system state well.During regular operation, a test pattern is declared as faulted if it falls outside this boundary and healthy in the contrary case (see [31,32] for further information about novelty detection).Due to its posttraining computational efficiency and ease of automation through a minimal number of configuration parameters, Tax's support vector data descriptor (SVDD) [33] for novelty detection is preferred by the authors.It provides many advantages over other traditional techniques [34].The SVDD fits target data with a minimal-radius hypersphere in an augmented space to generate nonconvex irregular decision boundaries in the normal feature space.The distance from the boundary is considered the novelty score; positive scores indicate that tested data fall within "normal behavior" while negative ones indicate a faulted state.
No Consideration of State.
When attempting to detect faults in a gearbox subject to varying load and speed, the impact of failing to consider the effects of these parameters can be severe.The results in Figure 4 demonstrate the consequences of using traditional fault-detection techniques that do not consider the variable nature of the problem; they are derived from a standard autoregressive model of order 20 and are fit to the vibration data that was in turn fit to an SVDD.While the healthy state is adequately characterized, all of the faulted states are so poorly indicated that it would be impossible to discern the presence of any of the described faults.The faults employed were relatively incipient in nature and one might assume that this approach might detect their presence later in the fault progression, possibly too close to catastrophic failure.
4.2.
Failure to Consider Load.The vibrations from a gearbox subject to both load and speed variations must be monitored with techniques sensitive to both parameters.Including the average speed of a feature vector's segment in that segment results in improved classification error and the earlier detection of faults as compared to those achieved when no efforts are made to adjust for time-varying parameters.Figure 5 demonstrates improved results that remain substantially poor.Severe faults like root cracks, chipped teeth, and outer race faults are easily detected due to the strength of their signals with respect to noise levels and the masking effects from speed and load variations.Less prominent faults like eccentric gears and more subtle bearing faults will remain masked without full consideration of all modal parameters.
State-Space
Based Feature Vectors.Taking full consideration of all predominant time-varying parameters drastically improves classification.Figure 6 demonstrates that all faults (subtle and severe) become easily discernible when employing state-space based feature vectors.The healthy class is somewhat difficult to classify having an error over 10%; this error is high and can be reduced via varying the order of the ARMAX models with a consequential tradeoff in classification error on faulted states.The analysis in the upcoming work exposes this approach's insensitivity to the double curse of dimensionality and its excellent tendency to generalize beyond untaught ranges of time-varying parameters [26].A more detailed discussion is limited herein but generalities are provided to facilitate a means of comparison.
Cross Correlation Model.
The cross correlation of vibration signals from disparate locations on a power train results in a signal not a feature vector; as discussed, this signal is in turn fit with an AR model whose parameters are used as the classification problem's feature vector.The vibration from the load motor's bearing was correlated with the vibration from the gearbox's input shaft bearing to generate the cross correlation results discussed.Figure 7 shows the effect of changing the model order on the classification results and the novelty score's distribution with respect to the decision boundary.Classification on all classes is poor with a low model order but as the model order increases the classification error drops in almost all cases.As was the case with state-space based feature vectors, a higher model order results in poorer classification error on the healthy class and good results on the faulted classes.This tradeoff seems present with cross correlation as there is a gentle increase in the error of the healthy state under these conditions.Balanced results are achieved with an order between 30 and 50 as shown in Figure 8.
Curse of Dimensionality.
The double curse of dimensionality arises when a large amount of data is not only required to characterize a high-dimensional system's behavior but when more is required due to the system's time-varying nature.Figure 9 demonstrates that this cross correlation technique enjoys a general immunity to the curse but it also demonstrates that classification results on the healthy training set can suffer with too little data.State-space based feature vectors were slightly less susceptible to this phenomenon [26].4.6.Generalization.The analysis surrounding Figure 9 and the double curse of dimensionality is relevant when analyzing these variable-state classification problems for generalization.This figure demonstrates that, with only limited training data from a select range of speed and load, cross correlation can be used to represent a gearbox's behavior in a manner not sensitive to these time-varying parameters.To analyze this effect further, consider classification results achieved when training is conducted with data from profiles shown in Figure 3 but with test data from Figure 10 (i.e., different accelerations on load and speed).Figures 11 and 12 demonstrate that the approach has excellent generalization when varied acceleration is used in speed and load; healthy and faulted data remain easily detected.
4.7.
Results from 80 : 48 Gearbox Arrangement.Figure 13 demonstrates that cross correlation works well with different mechanical parameters, that is, a less drastic gear ratio of 80 : 48.The gear teeth faults in this data set are fairly severe with concurring results in the form of novelty scores spaced a far distance from the decision boundary.4.8.Sensitivity Analysis: Segmentation Interval.Figure 14 demonstrates the effects of the length of the segmentation interval defined by a fixed number of (input) shaft rotations.There is a general trend of reduction in the classification error as the segmentation interval increases, particularly for the eccentric-gear fault; for most classes of fault, however, the change is not as dramatic.Classification is generally poor while the number of input shaft rotations falls below the gear ratio but becomes desirable after the interval rises to above 3-5 times the gear ratio.This seems reasonable as the output shaft will not have undergone a complete revolution until the former condition is met; after the latter condition is met, sufficient data to characterize the system's variations in a noisy environment will have been captured.Part (b) of this figure exposes a fairly consistent level novelty score distribution, suggesting that the improvement in classification error is not from a reduction in its variance but, instead, by a change in the average distance from the novelty boundary.
Automation
The block diagram in Figure 15 revises the steps taken in the proposed methodology as would be required in a more practical application.The segmentation interval could first be set to a value of 5 times the gear ratio.Vibration data would then be collected over these segments.The appropriate choice of the AR model order might vary from application to application; an online means of determining the appropriate order is therefore desirable.Since AR models are built by a least-squares fit of a signal's data samples, the choice of order could be selected by iterating through possible choices of order and using the one with the smallest 2 value.
Training data could then be collected over a number of segments fit to AR models which would in turn be stored until a certain amount had been collected at which point the SVDD boundary would be calculated.The completion of the SVDD training would be followed by online monitoring of the gearbox under scrutiny.
Conclusions
By cross correlating the signal from vibration sensors on disparate locations of a power train and processing the resultant signal into a feature vector for novelty detection, a powerful technique for classifying time-varying classification problems like fault detection in variable load and speed gearboxes has been demonstrated.The technique removes the need to analyze the complex nonlinear dynamics of the problem.It eliminates the need for costly sensors, like inline torque sensors, and the difficulties in deploying them in machinery not originally fitted for their use.The approach is computationally efficient and retains the excellent faultdetection abilities of other techniques under review.It also generalizes well across untrained state parameters.Through an established technique in sensor validation, the approach has been shown to provide a powerful means of reducing a complex condition monitoring problem to a near-trivial one.
Figure 2 :
Figure 2: Drive train (rigid shaft coupling (top left) turning pinion shaft with 32 teeth driving a gear with 96 teeth, rigid shaft coupling in line with torque meter supported by bungee cord, flexible coupling, and drive motor (bottom right)).
Figure 7 :
Figure 7: A sensitivity analysis of the AR model order fit to the cross correlation signal.
Figure 9 :Figure 10 :Figure 11 :Figure 12 :Figure 13 :Figure 14 :
Figure 9: The effects of the curse of dimensionality demonstrated by varying the amount of training data available.
Figure 15 :
Figure 15: Algorithm block diagram for practical online application. | 6,449.8 | 2014-12-22T00:00:00.000 | [
"Engineering"
] |
FEEDBACK NECESSARY OPTIMALITY CONDITIONS FOR A CLASS OF TERMINALLY CONSTRAINED STATE-LINEAR VARIATIONAL PROBLEMS INSPIRED BY IMPULSIVE CONTROL
. We consider a class of rightpoint-constrained state-linear (but non convex) optimal control problems, which takes its origin in the impulsive con- trol framework. The main issue is a strengthening of the Pontryagin Maximum Principle for the addressed problem. Towards this goal, we adapt the approach, based on feedback control variations due to V.A. Dykhta [4, 5, 6, 7]. Our necessary optimality condition, named the feedback maximum principle, is expressed completely in terms of the classical Maximum Principle, but is shown to discard non-optimal extrema. As a connected result, we derive a certain form of duality for the considered problem, and propose the dual version of the proved necessary optimality condition.
1. Introduction. In the paper, we address a particular class of state-linear optimal control problems with a simplest form rightpoint condition. This class of ordinary problems draws from an impulsive trajectory extension of a dynamical system, which is affine with respect to (w.r.t.) an unbounded control input: Here, T . = [0, T ] is a given finite control period, x 0 ∈ R n is an initial state position, U ⊂ R m is a given compact set, and A, B : U → R n×n , a, b : U → R n are given matrix and vector functions. System (1) is driven in two principally different ways: by an ordinary compact-valued measurable control u, and an L 1 -control v with unbounded values.
Integral relation in (2) is a sort of constraints that are sometimes called "soft" or "energetic bounds" [10], and makes sense of an a priori estimation (qualitatively expressed by a given positive real M ) of the total resource of controller over the time period T .
Systems of type (1), (2) and related control problems arise in a series of real life applications of mathematical control theory in mechanics, economics and management (see, e.g., the monographs [8,14,22], and the bibliography therein).
In fact, when regarding conditions (1), (2), one can immediately note that the model is ill-posed, since the control inputs v can be arbitrarily close in L 1 to distributions of Dirac type, and therefore, the respective state trajectories may tend to discontinuous functions. This assumes that related optimization problems could lack the existence of solutions, and therefore, are not correctly stated. Such an incorrectness can be overcome by a compactification of the trajectory tube of (1) in a relevant weak topology, weaker than the natural topology of uniform convergence (this could be the topology of pointwise convergence [2], or weak* topology of the Banach space of functions with bounded variation [14], or the topology defined by convergence of trajectories' graphs in the sense of Hausdorff distance [18]). The compactification leads to an impulsive dynamical system, i.e., a system with, possibly, discontinuous trajectories of uniformly bounded variation. Under certain natural convexity assumptions, the extended model takes the form of a dynamical system driven by distributions or measures: Here, x(0 − ) is the left one-sided limit of function x at zero (we agree that the trajectories of (3) are right-continuous), µ is a signed Borel measure, and |µ| denotes the total variation of µ. In fact, one should note that (3) is admitted to contain the product of measurable (discontinuous) function and a point-mass distribution, which is an incorrect operation. In this case, conditions (3), (4) are rather formal, though still give an intuition about the nature of the discussed trajectory extension. In general, the "control → trajectory" mapping of (3), (4) is set-valued, and a correct form of the desired trajectory extension is a more sophisticated measuredriven system with extra controls describing the way of approximation of Dirac type measures by ordinary controls. We refer to [1,2,14,15] for further details on impulsive trajectory extensions of dynamical systems, and pass to the wellknown but notable fact stating that the extended (relaxed) model (1), (2) can be equivalently transformed to an ordinary control system acting on the time interval S .
The cornerstone of this transformation is an absolutely continuous parametrization s → t(s) of the original time variable t = t(s) by a new "extended" time s. The original and reduced states are related by the (in general, discontinuous) time change x(t) = x(t ← (t)), where the symbol ← denotes the pseudo-inverse function defined by the relations: t ← (t) = inf{s ∈ S : t(s) > t} for t ∈ [0, T ) and t ← (T ) = T + M . This approach was independently invented by [17] and [21], and comprehensively treated in [14]. We note that system (5)-(7) has a specific dynamical structure, is weighted by a particular form terminal constraint, and is not convex in general. Based on the given notes, one can regard systems of type (5)- (7) and related variational problems as a substantial mathematical object, independently of their impulsive prototypes. The object of our present investigation is, thus, an optimal control problem for ordinary control systems of structure (5)- (7), and the main issue is a constructive strengthening of the Maximum Principle by virtue of technique [4,5,6,7], based on certain "feedback control variations". This strengthening is given by necessary optimality conditions of relatively new type. The conditions are extensions of the so-called "feedback minimum principle" to the discussed terminally constrained dynamics. The background of feedback optimality conditions is the technique of modified Lagrangians majorating the cost increment (the increment of an objective function), which -for state-linear problems -leads to an exact increment formula. These majorants are designed with the use of auxiliary functions, that are weakly monotone w.r.t. the system's dynamics [3]. By to the known principle of extremal aiming, such weakly monotone functions produce specific feedback controls, which can be used to define ordinary control processes that potentially "improve" a given reference solution. A wonderful fact is that feedback necessary optimality conditions can practically discard nonoptimal extrema of the Maximum Principle, and furthermore, can be thought of as iterative algorithms for optimal control.
2. Model statement. This section contains the setup of the main object of our study. Given a finite interval T = [0, T ], a compact set U ⊂ R m ; continuous matrix and vector functions A, B : U → R n×n , and a, b : U → R n ; vectors c, x 0 ∈ R n , and a positive real y T , consider the following optimal control problem (P ): As usual, ·, · denotes the scalar product in R n . A collection σ . = (z, w) .
= (x, y, u, v) is said to be a control process of system (8), (9), where 1], and • trajectories are absolutely continuous functions z .
A process σ is called admissible if it satisfies all the conditions (8)-(10).
3. Classical and feedback maximum principles. Letσ = (z,w) denote an admissible reference process, whose optimality is the question of our interest. Let us stress that problem (P ) is not assumed to be convex. It means that the Maximum Principle [16] does not turn here into a sufficient optimality condition, and therefore can be, potentially, improved. A rather challenging goal is to strengthen the classical result by earning an extra piece of information immediately from its standard relations. As we will show just below, such a strengthening can 204 STEPAN SOROKIN AND MAXIM STARITSYN be provided by feedback controls of a specific "extremal" structure, that remains in the formalism of the Maximum Principle.
Finally, note that the set of Carathéodory feedback solutions can be empty, while a sampling solution of (8), (9) does always exist. We denote Z(w) the set of solutions of both types (a) and (b).
3.2.
Formalism of the Maximum Principle. Introduce some necessary objects related to the Maximum Principle for problem (P ): The Pontryagin function (the non-maximized Hamiltonian) is written as are the "partial Hamiltonians" (note that H is independent of y); the adjoint (dual) system takes the form: The "variable" ξ = const is dual of y; note that, by the Maximum Principle, for extremal processes, ξ is not defined by a transversality condition. For processes, that do not satisfy the Maximum Principle, ξ can be thought of as a free parameter.
In what follows, this property will be essentially used in the definition of auxiliary feedback controls, which potentially discard local extrema. Denote Then, the maximized Hamiltonian takes the form: The maximizers are the following extremal multifunctions: (Here and on, Sign is the multivalued signature with Sign 0 = {−1, 1}.) Notice that the Maximum principle for the reference processσ = (z,w), in fact, reduces to the existence of an adjoint solution (ψ,ξ) such that the following inclusions hold a.e. on T : The idea of [4,5,6,7] consists in employing feedback controls w defined by a formal release of state position in inclusions (13): , v(t, x) ∈ Vξ(x,ψ(t)).
Terminal constraint.
It is clear that a solution of (8), (9) in any sense does not, generically, enjoy the condition y(T ) = y T . To take into account the terminal constraint, we introduce the following "corrected" multifunctions (in their definition, we employ the obvious description of the controllability set of trajectory component y to the point (T, y T )): Let W ξ denote the ξ-parametric, ξ ∈ R, set of feedback controls w = (u, v), which are selections of multivalued maps (14), (15) contracted to the dualψ of the reference trajectoryx, i.e., u(t, z) ∈Ǔ ξ (t, z,ψ(t)), and v(t, z) ∈V ξ (t, z,ψ(t)).
3.4.
Feedback maximum principle. Introduce the following accessory variational problem (AP ): The assertion below is a trivial implication of inclusions (13): Lemma 3.1. Letσ = (z,w) be a Pontryagin extremal for (P ). Then,z is admissible for (AP ), i.e., there exist ξ ∈ R and w ∈ W ξ such thatz coincides with a Carathéodory feedback solution z ∈ Z(w).
Proof. First, note thatσ is admissible for (AP ) by Lemma 3.1. Assume thatσ is optimal for (P ), but there exists ξ ∈ R, a feedback w ∈ W ξ and a respective feedback solution z = (x, y) ∈ Z(w) such that c, x(T ) < c,x(T ) . If z is a Carathéodory feedback solution, the contradiction is obvious. Let z be a sampling solution. Then, by its definition, there is a sequence z ρ of polygonal arcs converging to z in the uniform norm as diam ρ → 0. Since z ρ is a Carathéodory solution of (8), (9) produced by a piecewise constant open-loop control w ρ satisfying (10), (z ρ , w ρ ) is a control process of (P ), which still can violate the terminal constraint y(T ) = y T within accuracy diam ρ, i.e. |y ρ (T ) − y T | ≤ diam ρ. By shifting the rightmost point of ρ within its final segment T ρ (note that the length of T ρ does not exceed diam ρ!) we can perturb z ρ such that the perturbed functionz ρ = (x ρ ,ỹ ρ ) does satisfyỹ ρ (T ) = y T . Clearly,z ρ is also a Carathéodory solution of (8), (9) under another inputw ρ satisfying (10) (in fact w ρ = w ρ beyond T ρ ). Given ε > 0, consider a partition ρ such that diam ρ < ε and x(T ) − x ρ (T ) ≤ ε. Since the control perturbation is provided on a set of measure not exceeding diam ρ, standard arguments based on the Gronwall's inequality ensure that the quantity x ρ (T ) −x ρ (T ) has order ε, i.e., there exists a constant K > 0 independent of ε such that x ρ (T ) −x ρ (T ) ≤ K ε. Therefore, x(T ) −x ρ (T ) ≤ (K + 1) ε. Now, let ε > 0 be chosen such that c,x(T ) − x(T ) > (K + 1) ε c . Then, by the Cauchy-Bunyakovsky-Schwarz inequality, c,x ρ (T ) < c,x(T ) . Thus, we have defined a (P )-admissible process of a smaller cost. This contradicts the optimality ofσ. 4. Discussion and example. 1) Though the necessary global optimality condition proposed by Theorem 3.2 is hard for direct verification, even for a given process σ (not to say about applying this condition for searching "suspicious" processes), the feedback maximum principle is rather efficient in its counter-positive version, i.e., as a sufficient condition for non-optimality of a reference process. Indeed, for discardingσ as a non-optimal process, one can try to find a real ξ ∈ R and a feedback control w ∈ W ξ producing a feedback solution z = (x, y) ∈ Z(w) with the property: c, x(T ) < c,x(T ) = I(σ) (in this case, the admissibility ofz = (x,ȳ) for (AP ) does not require validation). Clearly, in such constructive form, the feedback maximum principle serves one as a conceptual kernel of iterative algorithms for problem (P ).
2) Note that the set of potentially discarding feedback controls and trajectories of (AP ) is extended due to the presence of the parameter ξ ∈ R. On the other hand, for practical implementation, the range of this parameter should be a priori estimated. In fact, an extremal processσ admits, generically, a number of adjoint trajectories (ψ,ξ) arising in Theorem 3.2. In general, one can identify the range of parameter ξ as follows: Assume we are given a multivalued map O(t) : T → R n , which contains the trajectory tube of system (8) started at (t, x) = (0, x 0 ). Then the parameter ξ is ranged in the interval [ξ − , ξ + ], where 3) The feedback maximum principle strengthen the classical Pontryagin Maximum Principle in the following sense: Assume that a functionz, being a Carathéodory solution of (8), (9) under a controlw satisfying (10), is admissible for (AP ). Then, σ = (z,w) is a Pontryagin extremal for (P ).
The absence of the inverse implication is shown by the following Example. Consider the following optimal control problem of type (P ): |v| ≤ 1.
STEPAN SOROKIN AND MAXIM STARITSYN
Finally, note that the considered model is a reduced version of the following optimal impulsive control problem:
It is clear that the cost functional of problem (P * ) can be reduced to the Mayer (terminal) form K(ψ, y, η, w) = η(0) + ψ(0), x 0 by introducing an extra state variable η satisfying the relationṡ After this transform, we can observe that the dual problem is also linear w.r.t. the state variable. Then, the developed feedback minimum principle can be adopted for problem (P * ).
The latter necessary optimality condition enjoys the same properties as the one proposed by Theorem 3.2. Note that, as certain examples show [7], similar conditions for free endpoint problems are not equivalent to each other, i.e. one of them does not imply the other one. 6. Conclusion. In the paper, we made a rather successful attempt to strengthen the Pontryagin Maximum Principle for a class of terminally-restricted variational problems. Though the derived optimality conditions -the feedback and dual feedback maximum principles -take variational forms, they can be though of as iterative algorithms for control improvement.
Note that, when implementing such algorithms, one inevitably faces similar issues in connection with the framework of discrete optimal control problems. In fact, similar results for discrete problems are recently obtained by the authors [19].
Finally, raising back the impulsive origin of the addressed specific class of models, one should point out the natural issue, which is left beyond this work. A rightful question here is a transcription of the obtained results into the terms of impulsive control. Such a transcription should exploit an appropriate correct notion of feedback impulsive control, which is not well introduced in the bibliography for measure-driven systems of the sort (3), (4). In this respect, the definition of feedback optimality conditions is a challenging issue of our nearest study. | 3,768.8 | 2017-06-01T00:00:00.000 | [
"Mathematics",
"Engineering"
] |
Effect of endogenous microbiota on the molecular composition of cloud water: a study by Fourier-transform ion cyclotron resonance mass spectrometry (FT-ICR MS)
A cloud water sample collected at the puy de Dôme observatory (PUY) has been incubated under dark conditions, with its endogenous microbiota at two different temperatures (5 and 15 °C), and the change in the molecular organic composition of this sample was analyzed by Fourier transform ion cyclotron resonance mass spectrometry (FT-ICR MS). Microorganisms were metabolically active and strongly modified the dissolved organic matter since they were able to form and consume many compounds. Using Venn diagrams, four fractions of compounds were identified: (1) compounds consumed by microbial activity; (2) compounds not transformed during incubation; (3) compounds resulting from dark chemistry (i.e., hydrolysis and Fenton reactions) and, finally, (4) compounds resulting from microbial metabolic activity. At 15 °C, microorganisms were able to consume 58% of the compounds initially present and produce 266 new compounds. For this cloud sample, the impact of dark chemistry was negligible. Decreasing the temperature to 5 °C led to the more efficient degradation of organic compounds (1716 compounds vs. 1094 at 15 °C) but with the less important production of new ones (173). These transformations were analyzed using a division into classes based on the O/C and H/C ratios: lipid-like compounds, aliphatic/peptide-like compounds, carboxylic-rich alicyclic molecule (CRAM)-like structures, carbohydrate-like compounds, unsaturated hydrocarbons, aromatic structures and highly oxygenated compounds (HOCs). Lipid-like, aliphatic/peptide-like and CRAMs-like compounds were the most impacted since they were consumed to maintain the microbial metabolism. On the contrary, the relative percentages of CRAMs and carbohydrates increased after incubation.
A fairly recent approach consists of infusing the whole preconcentrated and desalted cloud DOM sample into the ionization source with an ultrahigh resolution Fourier transform ion cyclotron mass spectrometer (FT-ICR MS). This nontargeted approach allows for the identification of the molecular formula C c H h N n O o S s in multiple compounds and allows for the computation of many useful parameters, such as elemental O/C (oxygen to carbon) and H/C (hydrogen to carbon) ratios 26 , the DBE (double bond equivalent) and the aromaticity index, which are useful parameters for estimating the carbon oxidation state, number of unsaturations and presence of aromatic compounds, respectively. Recent studies have used this powerful approach to investigate the molecular composition of cloud waters [26][27][28] , revealing the high degree of molecular complexity of this medium, with compounds related to both anthropogenic and biogenic sources in different oxidation states. For example, Cook et al. showed the influence of biogenic, urban and wildfire emissions on the molecular composition of cloud water samples collected at Whiteface Mountain (US) 28 , while Bianco et al. found a large contribution of biologically derived materials, such as lipids, peptides and carbohydrates, in cloud waters sampled at the puy de Dôme (PUY) station (France) 27 .
The main objectives of this study were (1) to evaluate variations in cloud DOM molecular diversity based on the activity of endogenous microflora in clouds and (2) to study how temperature can modulate this biological response. For this study, cloud water collected at PUY was incubated at two temperatures and analyzed by FT-ICR MS. Cloud processing is known to contribute to secondary organic aerosol (SOA) generation 10,25,29 , and this work will help to evaluate the impact of microorganisms on cloud organic matter transformations.
Results
Cloud water was collected on June 1 st , 2016, at the PUY station, between 2:50 pm and 7:20 pm. The air mass origin and FT-ICR MS analyses of this sample were discussed in a previous work 27 . This study showed that the puy de Dôme summit was located in the free troposphere during sampling, as indicated by the LACYTRAJ back-trajectory model. The sample was classified as marine using a multicomponent statistical analysis of the aqueous inorganic composition, as presented in Deguillaume et al. 20 . Continental influence may, however, not be excluded because the air mass was located below the boundary layer before arriving to the puy de Dôme station. Chemicophysical and microbiological characterization are reported in Supplementary Table S1. This table also reports the relative degradation or production of targeted chemical compounds in BIO + CHEM 5 and 15. The total number of cloud microorganisms was 8.6 × 10 4 cells mL −1 , as measured by the flow cytometry and microbial activity, which were confirmed by the ATP measurement (2.24 pmol mL −1 ). These values show that the microbial characteristics of this cloud sample are consistent with those in previous studies conducted on cloud water samples 7 .
The cloud water sample was filtered, preconcentrated on a solid phase extraction (SPE) cartridge and analyzed by FT-ICR MS. This fraction is referred to as "INITIAL" in the rest of the article. More detailed information about the analysis is given in the materials and methods section and in Supplementary S (1). Incubation experiments using the endogenous microbiota were then performed in the lab. Supplementary Fig. S1 summarizes all of the incubation tests that have been performed in this study. Cloud water was incubated in the dark, which allowed us to investigate the effects of dark chemistry (i.e., hydrolysis and Fenton reactions) and microbial metabolism. Two fractions containing endogenous microbial populations were incubated at 15 °C (BIO + CHEM 15) and 5 °C (BIO + CHEM 5) and preconcentrated and analyzed by FT-ICR MS with the same procedure used for INITIAL. Values of 5 °C and 15 °C were chosen as representative temperatures measured at the PUY station in winter and summertime, respectively. Transformations of the DOM values observed for these two fractions are attributed to both microbial metabolism and dark chemistry. To evaluate the contribution of dark chemistry only, two fractions of cloud water (CHEM 15 and CHEM 5) were also filtered to eliminate microorganisms, incubated in the same conditions as those used for BIO + CHEM 15 and 5 and analyzed by FT-ICR MS using the same set of experimental parameters. By comparing the compositions of BIO + CHEM and CHEM, the transformations resulting from microbiological processes only are highlighted. In terms of the acquisition of sequential mass windows, relative abundance was not considered in this study. In the first approach, the formation and degradation of organic compounds were evaluated on the basis of the comparison/disappearance of peaks (and associated assigned molecular formulas). Figure 1a shows van Krevelen (VK) diagrams corresponding to INITIAL, BIO + CHEM 15 and CHEM 15. Incubation at the highest temperature causes the degradation of a large fraction of compounds but also the formation of many new compounds. Figure 1a suggests the negligible impact of dark chemistry on the molecular composition of the sample. Indeed, the CHEM 15 (in blue in Fig. 1a) experiment did not result in any compound degradation, and only 25 compounds were produced when compared to the INITIAL sample. Since the dark chemistry contribution is non-significant, the result of the BIO + CHEM 15 incubation mainly reflects endogenous microbial metabolism. This incubation (BIO + CHEM 15) results in a significant decrease (42%) in the total number of compounds The analysis of the BIO + CHEM 5 and CHEM 5 results, compared with those of the BIO + CHEM 15 and CHEM 15 fractions, demonstrates the influence of temperature on endogenous cloud microbial activity. As observed for incubation at 15 °C, many compounds are degraded and formed, but some compounds are specifically produced at 5 °C. The effect of dark chemistry on molecular composition is minor: only 27 compounds are formed, which are shown in blue in Fig. 1b. The incubation of microorganisms at 5 °C leads to the production of a lower number of new compounds compared with the incubation experiment at 15 °C. Indeed, BIO + CHEM 5 contains 374 compounds, corresponding to approximately 20% of the initial number of assigned molecular formulas. In addition, produced compounds are different from the compounds produced during incubation at 15 °C: 190 compounds are common, and 173 are specifically formed at 5 °C. The average H/C values decrease by 7% at 5 °C, and this result is not observed at 15 °C ( Supplementary Fig. S4). This could be explained by the presence of 35 compounds with O/C ≤ 0.5 and H/C ≤ 1 in the incubated fraction (9.4% of the total assigned formula), while INITIAL contains only 9 compounds (0.5%). The yellow area represents compounds consumed by microorganisms (CONSUMED); the green area (NOT IMPACTED) corresponds to compounds not modified by chemical and biological processes; the blue area displays compounds produced by dark chemistry reactions (CHEM-PRODUCED); the pink area represents compounds produced by microorganisms (BIO-PRODUCED). The VK diagrams corresponding to each colored area of the Venn diagram are displayed below using the same color code.
Van Krevelen diagrams: effects of microbial and chemical transformations.
www.nature.com/scientificreports www.nature.com/scientificreports/ CONSUMED, NOT-IMPACTED, CHEM-PRODUCED and BIO-PRODUCED compounds, highlighting their distribution among the major biochemical classes (i.e., lipids, proteins, lignins, carbohydrates, and condensed aromatics) (Supplementary Table S2). Figure 3 reports the average number of carbon (nC), nitrogen (nN), hydrogen (nH), oxygen (nO) and sulfur (nS) atoms for CONSUMED, NOT-IMPACTED, CHEM-PRODUCED and BIO-PRODUCED. The NOT-IMPACTED compounds have significantly lower nC, nH and mass weight in comparison to the CONSUMED, CHEM-PRODUCED and BIO-PRODUCED compounds. The focus of this work deals with the effects of microbial degradation on organic matter; thus, the comparison between CONSUMED and BIO-PRODUCED is further analyzed. The carbon number (nC) in BIO-PRODUCED is higher than that in CONSUMED, and the values are more dispersed; the average nH does not significantly vary, while the nN and nS values increase by more than 10% compared to the CONSUMED values. The oxygen number (nO) decreases slightly with the weakening of O/C during incubation. (7) highly oxygenated compound (HOC) structures. CRAMs consist of carboxylated alicyclic structures and contain structures that are similar to large, fused, non-aromatic rings, with a high ratio of substituted carboxyl groups. Peptides are biomolecules consisting of short chains of amino acid residues. Lipid-like materials include fats, waxes, sterols, fat-soluble vitamins, monoglycerides, diglycerides and triglycerides. All of these molecular families are related to biological activities. Carbohydrate-like materials are complex polymeric structures composed of acyl polysaccharides, which contain varying levels of carbohydrate, lipid and acetate groups. HOCs are organo-nitrates and nitro-oxy organosulfates previously detected in aerosols and fog water samples [32][33][34][35][36][37] .
Effect of microbial activity: comparison of consumed and bio-produced.
A detailed composition of INITIAL has been previously described and discussed in Bianco et al. 27 and compared with other samples collected at PUY and Stormpeak Laboratory (US) by Zhao et al. 26 . Compounds contained in CONSUMED, BIO-PRODUCED, CHEM-PRODUCED and NOT-IMPACTED were analyzed following the same approach. Figure 4 presents the relative percentages of each class in the number of compounds. All of the classes are impacted by microbial transformations, with either an increase or decrease in the number of compounds. Endogenous microbiota preferentially degrade lipid-like and aliphatic/peptide-like materials. For lipid-like materials, the number of compounds decreases drastically from 303 to 44 in the assigned molecular formula; furthermore, the average weight of the molecular compounds in this class increases significantly after incubation (average values of CONSUMED = 398 ± 79 Da and BIO-PRODUCED = 540 ± 129 Da) (Supplementary Table S3). The number of molecular formulas assigned to the aliphatic/peptide-like class also strongly decreases from 430 (39.3% of the total) to 38 (14.3%). The CRAM-like compounds are produced during incubation. Surprisingly, incubation leads to the formation of unsaturated hydrocarbons (from 8 to 39 molecular assigned formulas) and aromatic compounds (formation of 2 molecules).
Comparison between incubations at 5 °C and 15 °C. A similar analysis was also performed for incubation at 5 °C. Figure 5 displays the Venn diagram and VK plot in the same way as Fig. 2. A comparison with the incubation at 15 °C reveals that a higher number of compounds is degraded (1716 at 5 °C instead of 1094 at 15 °C), but a lower number is produced (173 at 5 °C instead of 266 at 15 °C). As observed for incubation at 15°C, the molecular assigned formulas are distributed among the seven classes of compounds found in natural organic matter, except for NOT-IMPACTED, where the compounds are clustered in lipid-like and aliphatic/peptide-like regions. In the following paragraph, the comparison between CONSUMED and BIO-PRODUCED is detailed.
BIO www.nature.com/scientificreports www.nature.com/scientificreports/ and H/C ≤ 0.67). The median value of the molecular weight is higher for BIO-PRODUCED compared to CONSUMED. The comparison of BIO-PRODUCED between 15 and 5 °C shows that the molecular weight after incubation increases more at 5 °C than at 15 °C (26% instead of 6% of the average value).
Supplementary Figure S6 displays the relative abundances of compounds contained in the seven classes of compounds found in the DOM for the incubation experiments conducted at 5 °C. A decrease of 10.5% is observed in the relative percentage of the lipid-like compounds, and the mass weight decreases from 387 ± 81 to 333 ± 60 Da (Supplementary Table S3). In the aliphatic/peptide-like class, the average molecular weight (from 397 ± 89 to 555 ± 96 Da) increases with the formation of a molecular formula with more than 20 carbon atoms. The molecular weight is also higher for CRAM-like compounds (increase of 157 Da). As observed for incubation at 15 °C, compounds belonging to unsaturated hydrocarbon and aromatic classes are produced during incubation; in particular, incubation at 5 °C leads to the production of 17 compounds. The yellow area represents compounds consumed by microorganisms (CONSUMED); the green area (NOT IMPACTED) corresponds to compounds not modified by chemical and biological processes; the blue area displays compounds produced by dark chemistry reactions (CHEM-PRODUCED); the pink area represents compounds produced by microorganisms (BIO-PRODUCED). The VK diagrams corresponding to each colored area of the Venn diagram are displayed below using the same color code.
Discussion
These experimental results clearly highlight that endogenous microbiota alter the molecular composition of DOM in cloud water through the degradation of many compounds and the production of many others, while transformations by dark chemistry are negligible. Our results are consistent with those in previous experiments performed on three different cloud samples collected at the same site, with similar chemical and biochemical compositions 7 . In this study, only very limited numbers of short chain carboxylic acids (acetate, formate, succinate, oxalate and malonate), formaldehyde and hydrogen peroxide were considered, and microbial activity was largely prevailing with respect to dark chemistry. In the present work, we extended the investigation to a wider number of compounds presenting much larger molecular weights and more diverse chemical composition. Several hundred organic compounds were completely degraded by cloud microorganisms (1716 at 5 °C, 1094 at 15 °C), and hundreds of new compounds were biosynthesized (173 at 5 °C, 266 at 15 °C). Temperature is a stress factor for microorganisms in the atmosphere: at 5 °C, cell division is slowed down and after incubation the number of cell is lower at 5 °C than at 15 °C. In contrast, at low temperature, microorganisms need more energy to maintain metabolical activity and to produce ATP. For this reason, they consume more carbon sources, consuming more compounds during incubation.
One limitation to the present study is that ESI-FT-ICR MS is not a quantitative method. Only mass peaks that completely disappear and newly formed peaks are considered. Thus, the impact of microbiological activity on cloud chemistry is definitely underestimated. Many compound concentrations are changing as a consequence of microbial transformations but without the formation/disappearance of mass peaks. This can be observed regarding the biodegradation of formaldehyde, formate and acetate, which are not fully consumed during incubation (Supplementary Table S1). Even if they are impacted by microorganisms, they are part of the NOT-IMPACTED group following the presented classification. Abiotic transformations by hydrolysis and Fenton processes also partially degrade some compounds, but no variation in the number of compounds is detected. Moreover, solid phase extraction leads to the loss of some low molecular weight compounds.
Microbial metabolism is very complex and includes numerous pathways involved in building up and breaking down cellular components. Microbial extracellular enzymes may transform large molecules into smaller compounds in cloud water. Small molecules, which are initially present in the aqueous cloud phase or produced by aqueous phase reactivity, can be transported through the microbial membrane. Inside a microbial cell, small molecules are used to produce energy, converted into other small molecules or used to synthetize larger molecules. These large molecules, such as DNA or proteins, can be either integrated into the biomass or excreted in the cloud medium. Intracellular microbial metabolism also produces small molecules, which can be exported out of cells. Moreover, cloud microbiota are a dynamic system, where microbial cells are constantly growing or dying. As a consequence, cell lysis can occur and release both large and small compounds in cloud water. The impact of microbial activity on the cloud chemical composition is hard to evaluate since it is able to produce and consume organic compounds. For this reason, the CONSUMED and BIO-PRODUCED compounds may not necessarily be correlated.
Considering incubation at 15 °C, the H/C and O/C ratios do not vary significantly between CONSUMED and BIO-PRODUCED, but microbial transformations lead to the production of some more oxidized and reduced compounds (Fig. 1a,b). The DBE value and number of carbon atoms increase after incubation. The median value of the molecular weight remains constant, but more dispersion is observed with the combined production of both higher and lower weighted compounds. This could be due to the synthesis of high weight molecular compounds from small molecules and excretion out of the microbial cell. For example, cloud microorganisms were shown to produce extrapolymeric substances 12 , biosurfactants 38 and siderophores 39 . This biological activity is also observed in other environments, such as surface waters 40 . As reported before, the presence of high weight molecular compounds could also be the result of cell lysis.
Cold shock is a stress condition for microorganisms, and microorganisms handle this stress by modulating their metabolism 41 . A recent study conducted on Pseudomonas syringae isolated from cloud samples at the PUY station showed the metabolic effects of a temperature change from 17 °C to 5 °C 42 . This bacterium belongs to a more frequent and major active group in clouds 43,44 . Pseudomonas syringae synthetized cryoprotectants (namely, trehalose, glucose, glycerol, carnitine and glutamate) as a metabolic response to this cold shock 42 . The lipid metabolism was also altered by changing the saturation level of the fatty acids to preserve the membrane fluidity of the bacterium. In addition, the carbohydrate metabolism was activated to produce more energy (i.e., higher amounts of ATP), and the amino-acid metabolism was modified together with the synthesis and consumption of short di-and tetra-peptides.
In this work, we also observed a number of important changes in the metabolism of the cloud microbiome when incubations were performed at 5 °C compared to 15 °C. For example, incubation at 5 °C leads to the production of a lower number of compounds, which is different from those observed during incubation at 15 °C. Although the DBE median values are quite similar for both temperatures, a net shift toward higher molecular weights and nC is observed at 5 °C, as shown by the increase in the 75 th percentile (Fig. 3 and Supplementary Fig. S5). For these reasons, even if a large number of compounds is degraded, many compounds are also biosynthesized. This DBE increase could partially be related to the synthesis of unsaturated lipids to modulate the membrane fluidity in response to a cold shock, which is in agreement with previous reports 42 .
All chemical classes undergo a decrease in the number of compounds, except for aromatics and unsaturated hydrocarbons, whose numbers increase at both incubation temperatures. In addition, the impact of microbial transformations is different for each class, as highlighted by the parameter R, which is calculated using Eq. E1 (shown in Fig. 6) for both sets of incubations at 5 °C and 15 °C: www.nature.com/scientificreports www.nature.com/scientificreports/ where a = n° compounds in CONSUMED for the considered class, b = n° compounds in BIO-PRODUCED for the considered class, c = total n° of compounds in CONSUMED and d = total n° of compounds in BIO-PRODUCED. For R > 0, compounds are consumed; for R < 0, they are produced. We can observe that even if the number of compounds decreases in each class, lipid-like and aliphatic/peptide-like materials are degraded, while in the other classes, compounds are produced. Compounds related to the lipid-like class are consumed in both incubations (Fig. 6). These compounds may preferentially be degraded because of the high yield of ATP given by their oxidation, but they can also be incorporated into membranes during new cell formation. It is likely that microorganisms multiply during the incubation time scale, as previously observed by Amato et al. 45 . Low molecular weights of lipid-like materials are likely rapidly transformed into short chain carboxylic acids and other oxidized short molecules with great variations in the H/C and O/C ratios. For organic compounds in the lipid-like class with a molecular weight higher than 450 Da, CONSUMED contains 44 assigned molecular formulas, with an average mass of 524 ± 99 Da, H/C = 1.72 ± 0.10 and O/C = 0.24 ± 0.03. BIO-PRODUCED contains 44 assigned molecular formulas, with an average mass weight of 540 ± 129 Da, H/C = 1.72 ± 0.16 and O/C = 0.24 ± 0.04. Microorganisms are not able to degrade these compounds but can oxidize them. The first step of oxidation is usually hydroxylation; this explains the increase in the average mass weight of 16 Da (i.e., one oxygen atom). Thus, high weight lipid-like materials are weakly oxidized, and the small variations in their H/C and O/C values leaves them in the lipid-like class. This could justify the apparent lack of variation in their H/C and O/C ratios. Aliphatic/peptide-like compounds are also used by microorganisms for protein synthesis or as a source of nitrogen. For this class, the average molecular weight decreases, and the compounds are mostly degraded by microorganisms (Fig. 6). Table S3 shows that this class of compounds is degraded more at lower temperature (758 molecular formulas assigned in CONSUMED at 5 °C and 430 at 15 °C). This is likely because microorganisms need more energy to maintain their metabolism at low temperatures. A previous study reported that both amino acids and peptide metabolisms were modified when a cloud bacterium was incubated at a low temperature 42 .
CRAM-like compounds are formed during incubation. They present structural features that are common with those of terpenoids (i.e., membrane constituents and secondary metabolites in a wide range of prokaryotic and eukaryotic organisms 46 ). Hertkorn et al. reported that CRAM-like compounds are the decomposition products of biomolecules, as indicated by the prevalence of carboxyl groups and the pattern of increasing oxidation with decreasing molecular size 47 . This could explain the production of newly oxidized and smaller compounds in this class. Supplementary Figs S7 and S8 compare O/C vs. molecular weight for CRAM-like compounds in CONSUMED and BIO-PRODUCED at 15 °C and 5 °C, respectively. CONSUMED CRAM-like compounds, with a molecular weight range of 300-400 Da, are degraded to (probably) give BIO-PRODUCED compounds with a molecular weight range of 200-300 Da and a higher O/C. The core of the BIO-PRODUCED assigned molecular formulas, with a molecular weight range of 500-650 Da (right of the plot in Supplementary Figs S7 and S8), may result from the partial oxidation of lignin residues 48,49 , which have high molecular weights and fall in the same area as CRAM-like compounds 50 . Lignin residues are difficult to be ionized using ESI, but biodegradation may introduce oxidation, making them easily detectable. Carbohydrate-like compounds are produced during www.nature.com/scientificreports www.nature.com/scientificreports/ incubation, and the average molecular weight decreases. This result is consistent with the usual biotransformation of carbohydrate polymers into oligomeric compounds, resulting in short chain molecules. As an example, cellulose is biodegraded in oligosaccharides, resulting in dimers or monomers, which can be uptaken by the cells. The effect of temperature is clear for carbohydrate-like compounds: even though they are produced, a fraction is consumed to produce ATP. This fraction is larger for the incubation at 5 °C (52 compounds) than for the incubation at 15 °C (12 compounds); this is consistent with the activation of the carbohydrate route and the increase in the ATP concentration observed at low temperatures 42 . Saccharides can also be used to produce intracellular trehalose, which is a well-known cryoprotectant 42 .
The formation of aromatics and unsaturated hydrocarbons is also observed (Fig. 6), and this increase is particularly highlighted by the fact that INITIAL does not contain (or contains only a few) molecular assigned formulas in these classes. The presence of aromatics could result from lignin biotransformation, for instance, while unsaturated hydrocarbons are consistent with the increase in unsaturated lipids.
This study clearly indicates that microbial activity is able to strongly modify organic matter in clouds. This leads to the general assumption that cloud microorganisms are able to modify the chemical properties of aerosol particles and, thus, impact atmospheric processes. These processes were shown to synthetize higher molecular weight compounds, such as exopolymeric substances or oligosaccharides, which can interact with water. The cloud microbiota may then affect the formation of cloud condensation nuclei, and surface-active molecules could also be produced or released into the cloud water by cell lysis. These molecules are able to reduce the surface tension level of aerosol particles and could, thus, enhance their ability to be activated into cloud droplets [51][52][53] . These molecules could also migrate towards the surface phase, reducing hygroscopic water uptake and perturbing the efficiency of the mass transfer between the gas and aqueous phases of clouds. Siderophores could also be produced by microorganisms for the uptake of iron required for their metabolism. Theses complex molecules are able to strongly complex iron in cloud water, changing its redox cycle and role in the cloud oxidative capacity 39,54 .
All of these microbial transformations have been shown to be modulated by temperature. Levels of oxidants, such as hydrogen peroxide, are also key parameters controlling the complex and highly variable cloud microbiome metabolism 8 . This demonstrates the necessity of performing similar high-resolution mass spectrometry investigations on other incubations of various cloud samples with contrasted chemical and biological compositions. Cloud waters will be collected at very different geographical sites, exposed to different environmental conditions. It will be also important to combine photo-and biodegradation processes during incubation experiments to evaluate their potential synergic effect occurring simultaneously in clouds. A targeted analysis will be performed in parallel on the selected organic compounds strongly transformed during the incubations to assess the mechanisms and their modulations by microorganism.
Experimental Materials and Methods
The PUY station belongs to European atmospheric survey networks: ACTRIS (Aerosols, Clouds, and Trace gases Research Infrastructure) and EMEP (the European Monitoring and Evaluation Program). The PUY observatory also belongs to the GAW (Global Atmosphere Watch) stations. A dynamic one-stage cloud water impactor (cut off diameter of approximately 7 µm 55,56 ) was used to sample the cloud droplets. Before cloud collection, the aluminum impactor was cleaned and sterilized by autoclaving. The sample was stored in sterilized bottles; a fraction of the cloud water was filtered using a 0.22 µm nylon filter within 10 min after sampling to eliminate particles and microorganisms. The hydrogen peroxide concentration and pH were determined a few minutes after sampling. A fraction was frozen on site and stored in appropriate vessels at −25 °C until these samples were analyzed to estimate the ion concentrations by ion chromatography (IC), total organic carbon (TOC) by a TOC Shimadzu analyzer and iron concentration by the spectrophotometric method 57 . More details about the physicochemical analysis are reported in Bianco et al. 27 . Microbiological analyses were performed on nonfiltered cloud samples. Cell counts were performed by flow cytometry (BD FacsCalibur, Becton Dickinson, Franklin Lakes, NJ) on 450 μL triplicates, which were added with 50 μL 5% glutaraldehyde (0.5% final concentration; Sigma-Aldrich G7651) stored for <1 week at 4 °C. For analysis, the samples were mixed with 1 vol. of 0.02 μm filtered Tris-EDTA, with a pH of 8.0 (40 mM Tris-Base; 1 mM EDTA; acetic acid to pH of 8.0), and stained with SYBRGreen I (Molecular Probes Inc., Eugene, OR) from a 100X solution. The counts were performed for 3 min (or 100,000 events) at a flow rate of ~80 μL min −1 (precisely further determined by weighting). The ATP concentration was determined using the following procedure: 400 µL of the sample was strongly mixed in a microtube with 400 µL of the extractant B/S from the ATP measurement kit used (ATP Biomass Kit HS, Biothema) and stored in a frozen state until further analysis. The ATP concentrations were determined by bioluminescence 58 , as reported by Amato et al. 3 .
The cloud water sample was treated and analyzed by FT-ICR MS, filtering the signal by S/N > 5. In this way, mass peaks with low intensity (lower than 1 × 10 6 ), which could be detected or not detected as a function of the complexity of the matrix, are not considered. A fraction of the cloud water (200 mL) was filtered with a 0.22 µm nylon filter, and an aliquot was immediately frozen (referred to as INITIAL); the aliquot was incubated at 15 °C (CHEM 15) and 5 °C (CHEM 5) under the same conditions described in the following paragraph. Two fractions of cloud water with an endogenous microbial population were incubated in a sterilized Erlenmeyer flask at 15 °C (BIO + CHEM 15) and 5 °C (BIO + CHEM 5), with 200 rpm shaking, under dark conditions for 60 hours. Incubation time was chosen to maximize microbial transformations of dissolved organic matter. The temperatures were selected because they are representative of winter (5 °C) and summer (15 °C) conditions observed at the PUY summit under cloudy conditions 3 .
Each fraction was filtered with a 0.22 µm nylon filter and then prepared for analysis by ultrahigh-resolution FT-ICR mass spectrometry by using the method described by Zhao et al. 26 and the previously detailed analysis of the same sample in Bianco et al. 27 and Supplementary (S1). The high-resolution mass spectrometry analysis was performed using an FT-ICR mass spectrometer (Bruker) equipped with an electrospray ionization (ESI, Bruker) source set in negative ionization mode. The FT-ICR mass spectra were processed using the Composer software www.nature.com/scientificreports www.nature.com/scientificreports/ (Sierra Analytics, Modesto, CA): internal recalibration was performed, a peak list of signals, with S/N > 5, was generated, and the molecular formulas were assigned using the search criteria C 1-70 H 1-140 N 0-4 O 1-25 S 1 . The criteria described by Koch and Dittmar 59 were applied to exclude formulas that do not occur abundantly in natural organic matter (NOM): DBE must have been an integer value, 0.2 ≤ H/C ≤ 2.4, O/C ≤ 1.0, N/C ≤ 0.5, S/C ≤ 0.2, 2 ≤ H ≤ (2 C + 2) and 1 < O ≤ (C + 2). All peaks were used without filtering by relative abundance because relative abundance was strongly dependent on the ionization capability and not only related to the concentration, especially in a complex matrix with sequential acquisition. | 7,294 | 2019-05-21T00:00:00.000 | [
"Chemistry"
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Strayfield calculation for micromagnetic simulations using true periodic boundary conditions
We present methods for calculating the strayfield in finite element and finite difference micromagnetic simulations using true periodic boundary conditions. In contrast to pseudo periodic boundary conditions, which are widely used in micromagnetic codes, the presented methods eliminate the shape anisotropy originating from the outer boundary. This is a crucial feature when studying the influence of the microstructure on the performance of composite materials, which is demonstrated by hysteresis calculations of soft magnetic structures that are operated in a closed magnetic loop configuration. The applied differential formulation is perfectly suited for the application of true periodic boundary conditions. The finite difference equations can be solved by a highly efficient Fast Fourier Transform method.
Micromagnetic simulation are often used for the characterization of magnetic materials with a certain microstructure. Since the magnetic samples are very large, only a small part of the material can be simulated. A naive truncation of the magnetic domain would lead to strong shape anisotropy originating from surface effects. Periodic boundary conditions (PBCs) allow to eliminate this influence of the surface by modeling periodic images of the primary supercell.
Most well-known micromagnetic finite difference simulation packages like OOMMF 1 , MuMax 32 , magnum. fd 3 , magnum.af 4 and Fidimag 5 calculate the demagnetization field without PBCs using an analytic expression of the demagnetization tensor of homogeneously magnetized cubes 6 combined with an efficient FFT method making use of the convolution theorem 7 . Since this method is based on an integral formulation of the magnetic strayfield equations, incorporating PBCs requires the summation over an infinite sum of periodic images. Solutions have been proposed for 1D and 2D problems 8,9 , however an extension to 3D is not possible since the occurring sums are not absolutely convergent. Using point-dipoles instead of finite magnetized cubes seems to overcome this limitation and allows true 3D periodic boundary conditions 10 . In contrast to true PBCs which require an infinite summation, some codes use pseudo PBCs where the summation is truncated after a finite number of periodic images 2,3 . Apart from the easier implementation those methods are well suited for systems of intermediate size, where a finite number of periodic images is sufficient to model the complete sample. This principle has even been applied to FEM simulation where it is called macro-geometry 11 . Since finite element calculations are based on the differential form of the strayfield equation (real) PBCs can be directly applied by providing a proper cell-connectivity which represents the periodic structure.
We propose the application of PBCs for 3D problems based on a differential form of the strayfield equations both for finite elements (FE) and finite differences (FD). We focus on the efficient strayfield calculation because it is the most time-consuming part of micromagnetic simulations. Due to the long-range interaction the strayfield is usually solved using integral formulations. Since we use a differential formulation, the discretization using FE or FD with PBCs is straightforward, however it requires the solution of a sparse system of equations. In case of FD the use of a Fourier space method allows direct inversion of the system and offers a significant speed up.
Strayfield calculation using PBCs
The magnetic strayfield h of a given magnetization m can be calculated by means of magnetostatic Maxwell's equations. Since the magnetic strayfield is curl-free a scalar potential formulation can be used: www.nature.com/scientificreports/ where u is the magnetic scalar potential and the magnetic field can be calculated as h = −∇u . Proper boundary conditions need to be defined in order to obtain a unique solution.
If the magnetization is localized in a magnetic region ⊂ R 3 the problem is called an open-boundary problem with u = O( 1 r ) as r → ∞ . Since the boundary condition are not known at the surface of the magnet, direct use of the differential formulation (1) would require the discretization of an (infinite) air domain outside of the magnet. Accurate and efficient methods for solving open-boundary problems are often based on a corresponding integral formulation. In finite-element micromagnetics, the strayfield is usually solved in the differential form (1) and a hybrid method by Fredkin and Koehler 12 is used in order to resolve the boundary conditions on the magnetic surface by means of the boundary element method. In finite-difference micromagnetics, the direct integration of the strayfield tensor 6 combined with an efficient FFT based convolution is commonly used.
When dealing with large periodic structures, only a fraction of the magnetic sample can be discretized. Assuming open-boundary conditions is no longer valid in this case. Instead, a periodic magnetization can be assumed (see Fig. 1) and thus PBCs can be used as proper boundary conditions for u.
Finite element discretization. The finite element formulation is based on the weak form of the differential equation (1) with proper test functions φ i . ∂� would be the domain boundary with a corresponding unit normal n . In case of PBCs those surface intergrals vanish, because there is no physical domain boundary. For the application of PBCs a mapping of boundary nodes to their periodic images needs to be provided in order to eliminate the corresponding degree of freedoms. We used finite element packages FEniCS 13 or firedrake 14 (via firedrake-periodicity), which offer capabilities to define PBCs.
One difficulty when dealing with PBCs in FEM is the creation of a periodic mesh. Sophisticated periodic grain structures can be created with neper (see for example Fig. 2). Compared with a finite difference discretization the finite element model provides a better geometry representation, but requires the solution of an (unstructured) sparse system of equations, which significantly increases execution time by some orders of magnitude.
Finite difference discretization. The finite difference method is based on a rectangular N x × N y × N z mesh. In order to use an efficient FFT method for the solution of the occurring system of equations, we assume that the mesh is equidistant. Thus the index set (i, j, k) is sufficient to identify each vertex of the mesh.
In the following, we use the convention that u and div (m) are defined on the grid vertices, whereas m and h are defined at the cell centers (see Fig. 3). In the case of PBCs, the number of grid points and cell centers are equal and the cell centers form a shifted grid. www.nature.com/scientificreports/ The chosen convention allows to use central differences and leads to the following discrete approximations of the continuous differential operators Using these sparse discrete operators allows to solve the discrete version of Eq. (1) Efficient implementation using a fourier space method. The discrete system (5) can be solved directly, as it has to be done in the finite element method. Due to the regular (and equidistant) grid, the FD system can also be solved in Fourier space, where all differential operators become algebraic and the system can be directly inverted. The potential u i,j,k can be represented by means of the corresponding Fourier-space potential ũ i,j,k using the Discrete Fourier Transform (DFT) www.nature.com/scientificreports/ where is the imaginary unit and the wave-vectors k l , k m , and k n are defined by A similar Ansatz is used for the other fields m , and h . When substituting the Fourier-space representation (6) into the definition of the discrete operators (4) the spatial indices i, j, k only occur within the exponent, which results in a simple multiplicative phase factor for all neighbouring cells. Simplifying all occurring prefactors finally yields One can see that the non-local terms within the operator lead to local pre-factors within Fourier-space. Using the fact that the Fourier basis functions are linearly independent of each other, allows to explicitly express the Fourier coefficients of the strayfield h l,m,n as a function of the Fourier coefficients of the magnetization m l,m,n : Note that evaluation of ũ 0,0,0 , which represents the constant part of the potential, would lead to a division by 0. However it can be set to zero since it has no influence on the magnetic field.
Due to the use of the FFT, the resulting algorithm is very efficient. Compared with the non-periodic FFT strayfield calculation, which is based on the integral formulation, the assembly and storage of the demagnetization tensor can be avoided. Since no zero-padding is necessary the system size gets even smaller in case of PBCs. As a further optimization one can use a real-FFT, since the input magnetization as well as the resulting strayfield is purely real-valued. This leads to a further speedup by a factor of 2, as well as to a reduced storage size. The presented method is well suited for parallel execution on modern GPUs.
The FD method has been implemented within the micromagnetic code magnum.af and can be used on CPU or GPU, respectively. The FE method utilizing true PBCs has been implemented using magnum.pi. Table 1 summarizes timings of the strayfield calculation using the presented methods.
Numerical experiments
The presented method is validated by comparison with analytical calculations. The trivial case of a homogeneously magnetized bulk material leads to zero demagnetization field according to Eq. ((1)) since div m = 0 everywhere.
A non-trivial solution can be found for an infinite number of infinitely extended thin-films with thickness d 1 and a spacing between each thin-film of d 0 . Each thin-film is magnetized perpendicular to the film plane (in −x direction) with a saturation magnetization of M s . Due to symmetry considerations the resulting field only points (6) Note that in the limit d 0 → ∞ one ends up with the well known result for a single infinite thin-film with a field −M s inside of the film and 0 outside. It is remarkable that while a single thin-film shows no external field, an infinite number of films do. The calculated fields and the corresponding potential is visualized in Fig. 4 and compared with simulation results using the presented FD method. The performance of the presented method should be further demonstrated by the calculation of the hysteresis loop of a soft-magnetic-composite (SMC) material. The material consists of isolated particles, with each particle itself consisting of several magnetic grains. The simulation is restricted to a primary cell containing only one magnetic particle, consisting of 3 × 3 × 3 grains with a size of 300 nm × 300 nm × 300 nm , as well as a nonmagnetic interparticle layer (see Fig. 5). PBCs are used to mimic interparticle interactions and to avoid surface effects. The width of the interparticle layer w gap is varied and its influence on the magnetic hysteresis is studied.
For the FE model a (non-equidistant) regular mesh is used. Each dimension is divided into N = 3N i + 1 parts, where N i is the number of divisions of each grain. The grid spacing within the grains is constant, whereas the thickness of the interparticle layer can be adjusted as desired. In contrast, the FD discretization requires using an equidistant grid, which limits the possible interparticle thicknesses to integer divisors of the grain-size. Furthermore many FFT libraries require that largest prime factor of the system size is smaller than a certain value. This is based on the fact, that the FFT performs best for system sizes N = 2 M for integer M. Performance www.nature.com/scientificreports/ decreases dramatically (more than one order of magnitude) if system sizes with much higher prime-factors are used 2 . This general limitation of the FD method also results in very large system sizes for small width of the inter-particle layer and makes the more flexible FE method still competitive. The used material parameters can be found in Table 2. The magnetic anisotropy axes within the 27 grains are randomly distributed. A homogeneous external magnetic field is applied and linearly varied from −100 mT to 100 mT with a frequency of 100 MHz . The resulting hysteresis loops for FE and FD method can be found in Fig. 6. It can be seen that a larger interparticle layer leads to a smaller hysteresis and thus reduces the hysteresis losses of the material. The dramatic influence of self-demagnetization without using true PBCs is demonstrated in Fig. 7. Due to the strong demagnetization effects, the external field range needs to be extended to −2.5 T to 2.5 T . Without PBCs the subtle effect of varying interparticle layer widths is superimposed by a much stronger finite size-effect, which additionally depends on the shape of the boundary. Since the influence of the boundary can be subtracted out only in average, extracting the desired macroscopic material properties will be much harder and less accurate.
The FE and FD simulations with true PBCs have been performed with magnum.pi, or magnum.af, respectively. The FD simulations without PBCs have been performed with magnum.af and validated with MuMax 3 . The FD simulations using pseudo PBCs have been performed with MuMax 3 .
Conclusion
The importance of using true PBCs for the calculation of material properties without the influence of surface effects has been pointed out. An efficient FFT-based FD strayfield calculation providing true 3D PBCs has been presented. This method perfectly complements methods for 1D and 2D periodic boundary conditions. Those Figure 5. Simplified geometry of the SMC material consisting of 3 × 3 × 3 magnetic grains separated with one non-magnetic interparticle layer (dark blue). The size of each grain is 300 nm × 300 nm × 300 nm , whereas the interparticle thickness is varied. Table 2. Micromagnetic material parameters used within the magnetic grains as well as for the non-magnetic interparticle layer. The easy axes of the uniaxial anisotropy are randomly distributed. www.nature.com/scientificreports/ Finite difference hysteresis curves of the SMC material for interparticle widths w gap = 23.08 nm with different boundary conditions. Without true PBCs the external field range has to be extended to ±2.5 T , in order to fully saturate the material. | 3,263.4 | 2021-04-28T00:00:00.000 | [
"Physics",
"Engineering"
] |
Search for new non-resonant phenomena in high-mass dilepton final states with the ATLAS detector
A search for new physics with non-resonant signals in dielectron and dimuon final states in the mass range above 2 TeV is presented. This is the first search for non-resonant signals in dilepton final states at the LHC to use a background estimate from the data. The data, corresponding to an integrated luminosity of 139 fb$^{-1}$, were recorded by the ATLAS experiment in proton-proton collisions at a centre-of-mass energy of $\sqrt{s} = 13$ TeV during Run 2 of the Large Hadron Collider. The benchmark signal signature is a two-quark and two-lepton contact interaction, which would enhance the dilepton event rate at the TeV mass scale. To model the contribution from background processes a functional form is fit to the dilepton invariant-mass spectra in data in a mass region below the region of interest. It is then extrapolated to a high-mass signal region to obtain the expected background there. No significant deviation from the expected background is observed in the data. Upper limits at 95% CL on the number of events and the visible cross-section times branching fraction for processes involving new physics are provided. Observed (expected) 95% CL lower limits on the contact interaction energy scale reach 35.8(37.6) TeV.
Introduction
Signatures with dilepton (dielectron and dimuon) final states have been central in shaping the Standard Model (SM) over many years, from discoveries of new particles [1][2][3][4][5], through many precision measurements [6][7][8][9], and in searches for new physics beyond the SM (BSM) [10][11][12][13]. This has been the case owing to the clean and fully reconstructable experimental signature with excellent detection efficiency. This paper presents a novel search for new phenomena in final states with two electrons or two muons in 139 fb −1 of data collected in proton-proton (pp) collisions at the LHC at a centre-of-mass energy √ s = 13 TeV between 2015 and 2018. The work presented here complements the ATLAS search for heavy resonances [10] using the same dataset and selection criteria. The new physics signature investigated is a broad, non-resonant excess of events over a smoothly falling dilepton invariant-mass spectrum, which is dominated by the Drell-Yan (DY) process. The search results in this paper are provided in a model-independent format. These results are further interpreted in the context of the frequently tested benchmark models with effective four-fermion 'contact' interactions (CI) [14,15].
A number of changes are introduced with respect to the previous ATLAS result with an integrated luminosity of 36.1 fb −1 [13]. The result presented here is the first non-resonant dilepton search at the LHC to use a background estimate from the data using a functional form. The signals considered are expected to manifest themselves only as a deviation from the expected gradient of the high-mass tail of the dilepton mass spectrum. Therefore, the background at high masses is estimated from a low-mass control region (CR) where the signal contribution is expected to be negligible. Contrary to previous ATLAS searches for non-resonant signals in dilepton final states, this search is performed in a single-bin high-mass signal region (SR). Both the function and region choices are optimised to maximise the expected sensitivity to observe CI processes. The extrapolated background is integrated in the SR to provide an estimate of the expected number of background events. The signal would be seen as an excess over this expected background estimate. This CR/SR approach is essential in the case of (typically small) non-resonant signals, as when the entire mass range is fit, similar to Ref. [10], a non-resonant signal can be absorbed into the background model. Moreover, the choice of a single-bin signal region removes the dependence on the shape of the mass distribution and simplifies the entire search, while at the same time providing model-independent results.
Further, this analysis has been moved from a Bayesian statistical framework to a frequentist statistical framework, which removes the dependence on signal priors. In the case where the interference between signal and SM processes is not negligible, e.g. for CI, the choice of one prior over another is less justified [13,16]. With respect to the previous study [13] that used simulation to estimate the background, the approach presented here reduces the dependence on simulation by estimating the background from the data. A comparison showed little difference in sensitivity between the two approaches.
Finally, the transition to a background estimation from the data exchanges the systematic uncertainties in the predictions from simulation for statistical uncertainties in data. The dominant uncertainty in the expected background in the new analysis is due to statistical fluctuations in the CR. Next in importance is the uncertainty in the degree to which the extrapolation from the CR can produce a background estimate different from the underlying distribution, leading to a signal-like deflection in the SR. This uncertainty is quantified using the simulated background and its uncertainties. The uncertainty third in importance is due to a possible signal contamination in the CR.
Contact Interactions
In the SM, it is assumed that quarks and leptons are fundamental point-like particles and hence have no structure. However, if quarks and leptons are composite, with at least one common constituent, the interaction of these constituents could manifest itself through an effective four-fermion contact interaction at energies well below the compositeness scale [14,15], Λ, the energy scale below which fermion constituents are bound. A broad class of CI models can be described by the CI Lagrangian of the form of Eq. (1): where g is a coupling constant chosen such that g 2 /4π = 1, γ are the 4 × 4 Dirac matrices and the spinors ψ L,R are the left-handed and right-handed fermion fields, respectively. The parameters η i j , where i and j are L or R, define the chiral structure (left or right) of the new interaction. Specific models are chosen by assigning the parameters to be −1, 0 or +1.
In the context of CI searches with dilepton final states at the LHC, the terms in Eq. (1) take the form of η i j q i γ µ q i ¯ j γ µ j , where q i and j are the quark and lepton fields, respectively. The differential cross-section for the process qq → + − , in the presence of CI, can be separated into the SM DY term plus terms involving the CI. This separation can be seen in Eq. (2): where the first term accounts for the DY process, the second term corresponds to the interference between the DY and CI processes, and the third term corresponds to the pure CI contribution. The latter two terms include F I and F C , respectively, which are functions of the differential cross-section with respect to m with no dependence on Λ [14]. The interference can be constructive or destructive and it is determined by the sign of η i j .
Previously, the ATLAS and CMS experiments have searched for CI with the partial Run 2 datasets at √ s = 13 TeV [12,13]. The most stringent exclusion limits for qq + − CI, in which all quark flavours contribute, come from the previous ATLAS non-resonant dilepton analysis conducted using 36 fb −1 at √ s = 13 TeV. The observed lower limits on Λ range from 24 to 40 TeV depending on the specific signal model [13].
ATLAS detector
ATLAS [17-19] is a multipurpose detector with a forward-backward symmetric cylindrical geometry with respect to the LHC beam axis.1 The innermost layers consist of tracking detectors in the pseudorapidity range |η| < 2.5. This inner detector (ID) is surrounded by a thin superconducting solenoid that provides a 2 T axial magnetic field. It is enclosed by the electromagnetic and hadronic calorimeters, which cover |η| < 4.9. The outermost layers of ATLAS consist of an external muon spectrometer (MS) with |η| < 2.7, incorporating three large toroidal magnetic assemblies with eight coils each. The field integral of the toroids ranges between 2.0 and 6.0 Tm for most of the acceptance. The MS includes precision tracking chambers and fast detectors for triggering. A two-level trigger system [20] selects events to be recorded at an average rate of 1 kHz. process and the programs used for parton showering are listed in Table 1 with their respective parton distribution functions (PDFs). 'Afterburner' generators such as P [24] for the final-state photon radiation (FSR) modelling, M S [25] to preserve top-quark spin correlations, and E G [26] for the modelling of cand b-hadron decays, are also included in the simulation. The DY [36] and diboson [37] samples were generated in slices of dilepton mass to increase the sample statistics in the high-mass region. Next-to-next-to-leading-order (NNLO) corrections in quantum chromodynamic (QCD) theory, and next-to-leading-order (NLO) corrections in electroweak (EW) theory, were calculated and applied to the DY events. The corrections were computed with VRAP v0.9 [38] and the CT14 NNLO PDF set [39] in the case of QCD effects, whereas they were computed with MCSANC [40] in the case of quantum electrodynamic effects due to initial-state radiation, interference between initial-and final-state radiation and Sudakov logarithm single-loop corrections. These are calculated as mass-dependent K-factors, and reweight simulated events before reconstruction. The top-quark samples [41] are normalised to the cross-sections calculated at NNLO in QCD including resummation of the next-to-next-to-leading logarithmic soft gluon terms as provided by T ++2.0 [42].
All fully simulated event samples include the effect of multiple pp interactions in the same or neighbouring bunch crossings. These effects are collectively referred to as pile-up. The simulation of pile-up collisions were performed with P v8.186 using the ATLAS A3 set of tuned parameters [43] and the NNPDF23LO PDF set, and weighted to reproduce the average number of pile-up interactions per bunch crossing observed in data. The generated events were passed through a full detector simulation [44] based on G 4 [45].
In order to reduce statistical uncertainties, a large additional DY sample is used where the detector response is modelled by smearing the dilepton invariant-mass with mass-dependent acceptance and efficiency corrections, instead of using the CPU-intensive G 4 simulation. The relative dilepton mass resolution used in the smearing procedure is defined as (m − m true )/m true , where m true is the generated dilepton mass at Born level before FSR. The mass resolution is parameterised as a sum of a Gaussian distribution, which describes the detector response, and a Crystal Ball function composed of a secondary Gaussian distribution with a power-law low-mass tail, which accounts for bremsstrahlung effects or for the effect of poor resolution in the muon momentum at high p T . The parameterisation of the relative dilepton mass resolution as a function of m true is determined by a fit of the function described above to simulated DY events at NLO. A similar procedure is used to produce a mass-smeared tt sample. These two samples replace the equivalent ones produced with the full detector simulation wherever applicable in the remainder of the analysis. The number of events in these samples is more than 55 times the number of events in data. These samples would have been difficult to produce with the full detector simulation because of the large number of events required and the limited computing resources.
Signal m distribution shapes are obtained by a matrix-element reweighting [13] of the leading-order (LO) DY samples generated in slices of dilepton mass. This reweighting includes the full interference between the non-resonant signal and the background DY process. The weight function is the ratio of the analytical matrix-elements of the full CI (including the DY component) and the DY process only, both at LO. It takes as an input the generated dilepton mass at Born level before FSR, the incoming quarks' flavour and the CI model parameters (Λ, chirality states and the interference structure). These weights are applied to the LO DY events to transform these into the CI signal shapes, in steps of 2 TeV between Λ = 12 TeV and Λ = 100 TeV. Dilepton mass-dependent higher-order QCD production corrections for the signals are computed with the same methodology as for the DY background, while electroweak corrections are applied in the CI reweighting along with the interference effects. These signal shapes are used for optimisations as well as for calculations of the cross-section and acceptance times efficiency.
The statistical analysis used in this work requires a continuous description of the CI signal shape between the fixed (reweighted) signal shapes, for the values of Λ mentioned above. A bin-by-bin morphing procedure is used to obtain a smooth description as a function of Λ, linearly interpolating between the fixed signal shapes from simulation. In the case of constructive interference the morphing is almost redundant since the signal behaviour between different Λ values can be approximated with a relatively simple relationship between the signal strength and Λ. However, in the case of destructive interference there is no straightforward relationship between the signal strength and Λ, and so the morphing approach is essential. The morphing is only performed for values of Λ inside the range of the reweighted signals described above.
Object reconstruction and event selection
A complete description of the object definition and event selection is given in Ref. [10]. These criteria are identical to the ones used in this work and a brief description follows. The dataset was collected during LHC Run 2 in stable beam conditions, with all detector systems operating normally and while fulfilling all quality requirements. Events in the dielectron channel were recorded using a dielectron trigger, while events in the dimuon channel were required to pass at least one of two single-muon triggers. Further, it is required that at least one pp interaction vertex be reconstructed in the event. The events are required to contain at least two same-flavour charged leptons consistent with the primary vertex. The object definitions, single-lepton selection and corrections are given in Ref. [10]. The reconstruction of the same energy deposits as multiple objects is resolved using overlap-removal procedures. If more than two leptons are present in the event, the two leptons with the largest E T (p T ) in the electron (muon) channel are selected to form the dilepton pair. In events with a dielectron pair and a dimuon pair, the dielectron pair is selected because of the better resolution and higher efficiency for electrons. A selected muon pair must contain oppositely charged muons. For an electron pair, the opposite-charge requirement is not applied because of the higher probability of charge misidentification for high-E T electrons. The reconstructed mass of the dilepton system after the full analysis selection, m , is required to be above 130 GeV to avoid the Z boson peak region, which cannot be described by the same parameterisation as the high-mass part of the dilepton distributions.
Background modelling
The dilepton invariant mass distribution in data is fit by a parametric background-model function in a low-mass control region (CR). The resulting background model is then extrapolated from the CR to higher-mass single-bin signal regions (SRs). The normalisation of the background model in the CR is determined by the number of data events in the CR only. All fits are performed within the RooFit [46] framework. Different choices of CR and SR are considered in order to maximise the expected sensitivity for each lepton channel and for different choices of the CI model parameters. In the destructive interference cases, if a SR includes a significant part of the destructive component of the signal shape, the integral of the number of expected signal events in the SR is reduced. Therefore, the optimisation procedure allows a gap between the CR and SR to avoid the cancellation due to the range where the signal contributes destructively. The final CR and SR choices are checked to ensure that the possible presence of a non-resonant signal does not bias the background estimation in the CR and consequently also in the SR. An illustration of the division into CR and SR is shown in Figure 1. The monotonically falling total background shape is shown by the solid black line, while an example of a CI signal plus the total background shape is shown by the dotted red line. This CI signal shape corresponds to the last two terms in Eq. (2) for a destructive interference case. The two axes in the figure are shown in logarithmic scale. The data is fit in a low-mass control region (shaded blue area) where a potential bias from the presence of a signal is negligible. The resulting background shape is extrapolated from the control region into the high-mass signal region (shaded red area). The gap illustrated between the CR and the SR is found to be the preferred case for the destructive interference cases only.
An optimisation procedure is performed in two consecutive steps. In the first step, the fit function is chosen out of about 50 initial functions, which are all checked in a set of about 15 potential CR and SR configurations. Once the function choice is fixed, the CR and SR choice is optimised in a second step using this function. The description of these two steps is given below.
The procedure to determine the functional form of the background is as follows. The smooth functional form used to model the background is chosen from about 50 candidate functions. Each function is fit to the dilepton mass background template, consisting of the sum of all the simulated background contributions, in a variety of CRs and extrapolated to the respective SRs. The fit to data and simulation are both performed with a bin width of 1 GeV. The distribution of the pulls, defined as (fit-simulation)/fit for each bin, is obtained for each potential configuration of CR and SR. A function that results in pulls below 3 across all the ranges considered (CRs and SRs) is marked as acceptable. This requirement is particularly important in the SRs to veto functions that exhibit unphysical behaviour at the tail. Additionally, it is important to ensure a good description of the simulated background template in the CRs. Out of about 50 initial functions, five are found to satisfy this requirement equally well. The residual mis-modelling by the selected function is measured later and taken as an uncertainty. The final function is chosen to be the same one used in Ref. [10] and it is given in Eq. (3): , is a non-relativistic Breit-Wigner function with m Z = 91.1876 GeV and Γ Z = 2.4952 GeV [47]. The second term, (1 − x c ) b , ensures that the background shape evaluates to zero at x → 1. The parameters b and c are fixed to values obtained from fits to the simulated background. In the third term, the parameters p i with i = 0, .., 3 are left free in the fits. The function f b (m ) is treated as a probability density function in the fits performed in the CR. This function is then normalised in the CR to N CR , the number of events in the CR in data (or simulation where applicable), where it is assumed that the CR is completely dominated by background events.
After fixing the function choice, the procedure to define the CR and SR is as follows. The two boundaries of the CR (CR min and CR max ) and the lower boundary of the SR (SR min ) are chosen to optimise the expected sensitivity for each of the CI signals considered. The CR min value is varied between 160 GeV (well above the Z peak) and 500 GeV, while CR max is varied between 1 TeV and 2.9 TeV. The CR is not wide enough to constrain the fit for CR max values below 1 TeV, while above 2.9 TeV the possible new signals contribute significantly. In all cases, the upper boundary of the SR is fixed to 6 TeV, beyond the highest-mass events expected in data, while SR min can lie at any point above CR max . The boundaries of the CR are varied to test for a possible dependence of the background estimation in the SR and it is found that the estimation remains stable against these variations.
To avoid a bias from possible signal contamination in the CR, the CR and SR choice is validated using a signal injection test for each of the configurations tested. The signals are injected in the range 18 TeV ≤ Λ ≤ 40 TeV. A collection of background+signal distributions are produced by simulation for various Λ values of interest. An extension of Eq. (3), with an added signal component, is used to fit these distributions: where f s (m , Λ) is the signal probability density function and N s (Λ) is the number of signal events in the CR. Both f s (m , Λ) and N s (Λ) are determined from simulation. The parameter N b is the number of background events in the CR with the constraint N b + N s (Λ) = N CR . The full shape is fitted in the CR using the background+signal model and compared with the nominal case, where there is no signal injected and where the fit model is the background-only one. If a significant difference is found between the background estimated with the injected signal fit and the nominal background-only fit, the configuration is excluded. The difference between these two background estimates is assessed to be significant when it is larger than the systematic uncertainty of the background component in the background+signal model. This procedure is repeated iteratively while varying two out of the three mass boundary parameters (CR min , CR max and SR min ). It is found that the background component of the background+signal model does not differ significantly from the simulated background, both in the presence and absence of an injected signal.
Each chirality choice of the CI model is tested with an independent CR and SR configuration. It is found that for models with destructive interference a mass gap between the CR and SR of ∼ 1300 GeV is preferred by the optimisation procedure, while in the case of constructive interference the optimal choice is where SR min coincides with CR max . The resulting ranges for the different chirality options are similar at the level of a few tens of GeV. The final result is insensitive to the choice of CR within these small differences, and therefore these are merged as listed in Table 2 to simplify the subsequent procedures. The function given in
Uncertainties
Uncertainties related to the background modelling in the SR result from three main sources as discussed below. For all background variations discussed, where the extrapolation procedure is performed, it is verified that the χ 2 /N DoF of the fits to each of the background variations in the CR is close to unity. The uncertainties related to the signal model are also presented.
Statistical uncertainty of the expected background
Statistical fluctuations in the data lead to variations of the fitted background model in the CR. This in turn has an impact on the extrapolated background in the SR. To estimate the impact of this statistical uncertainty, σ Stat b , the following procedure is performed for each region configuration. First, the data is fit in the CR, extrapolated, and integrated in the SR, giving the nominal background expectation in the SR. The nominal background m distribution shape in the CR is then used as a probability density function from which an ensemble of pseudo-datasets can be generated. The normalisation of this function corresponds to that of the observed data in the CR. Finally, the background model is fit to each of the pseudo-datasets in the ensemble individually, extrapolated, and integrated in the SR. The distribution of the pseudo-background expectations is confirmed to be centred around the nominal background expectation, indicating no bias. The standard deviation of the distribution is taken as the statistical uncertainty. For the dielectron and dimuon channels, the statistical uncertainty ranges from 14% to 20% (34% to 60%) of the nominal background for the constructive (destructive) SRs of the analysis.
Induced spurious-signal uncertainty in the expected background
The second uncertainty in the expected background corresponds to the degree to which the background model can induce a signal-like excess or deficit when extrapolated to the SR. This uncertainty is hereafter called 'induced spurious-signal' (σ ISS b ). This uncertainty results from the extrapolation procedure and it is measured on the nominal simulated background and its systematic variations.
The uncertainties associated with the simulated background shape are derived from simulated variations on the background shape. These uncertainties are used to generate an ensemble of possible (pseudo-) background shapes. Each pseudo-background shape is constructed from the nominal simulated background shape, summed with weighted uncertainties. The weight for each uncertainty is randomly sampled from a normal distribution, with a mean of zero and standard deviation of one, in the range of [−1, 1]. The resulting background shape is used to generate a pseudo-dataset that is fit and extrapolated to the SR.
The difference in expected background between the fit and the pseudo-background SR integral is then taken as the induced spurious-signal per pseudo-background. The mean and standard deviation of the distribution from all pseudo-backgrounds are summed in quadrature and the result is taken as σ ISS b . The mean is considered to take into account a possible systematic shift in the estimate besides its spread.
The variations considered are due to theoretical and experimental uncertainties in the simulated background as well as the uncertainties in the backgrounds from multi-jet and W+jets processes. The largest source of uncertainty in the simulated background is theoretical, and it is particularly large at the high end of the dilepton mass spectrum. The second largest source of uncertainty in the simulated background is experimental, and is mostly due to high-p T muon identification in the dimuon channel. The third largest source is the uncertainty in the multi-jet and W+jets background components, and is estimated from the data.
The following variations are considered for the theoretical uncertainties for the DY component only: the eigenvector variations of the nominal PDF set, variations of PDF scales, the strong coupling (α S (M Z )), electroweak corrections, photon-induced corrections [48], as well as the effect of choosing different PDF sets. For all PDF variations, the modified DY component is used along with the other nominal background components. These theoretical uncertainties are the same for both dilepton channels at generator level, but they result in different uncertainties at reconstruction level due to the different resolutions of the dielectron and dimuon channels. Further details of this procedure can be found in Ref. [13]. The size of these uncertainties in the total simulated background is ≤ 19% (≤ 15%) below 4000 GeV for the dielectron (dimuon) channel.
Among the experimental uncertainty sources in the dielectron channel, the dominant ones are the electron identification at low dielectron masses (≤ 5%, below ∼ 2000 GeV) and the uncertainty in the electromagnetic energy scale at higher dielectron masses (≤ 15%). In the muon channel, the dominant experimental uncertainties arise from the muon reconstruction efficiency at low dimuon masses (≤ 20%, below ∼ 4000 GeV) and from the identification of high-p T muons at higher dimuon masses (≤ 50%).
The relative uncertainty of the simulated background due to the multi-jet and W+jets component rises from ∼ 1% at 1 TeV to ∼ 10% at 4 TeV. For the multi-jet and W+jets component variations, the modified shape is used each time along with the other nominal background components from simulation. This contribution is the smallest amongst all other variations in the CR.
The σ ISS b uncertainty is ∼ 4% (∼ 6%) for the constructive e + e − (µ + µ − ) channels, and is ∼ 7% (∼ 24%) for the destructive channels. The large difference between the µ + µ − and e + e − channels in the destructive case is owing to the smaller CR in the µ + µ − case as can be seen in Table 2. Consequently, the µ + µ − background fit in the CR is less constrained, allowing for more freedom in the extrapolation to the SR.
CR bias uncertainty in the expected background
Finally, the 'CR bias uncertainty' (σ CRB b ) in the expected background is a measure of the residual difference between the two fit models, with and without a signal component. A possible signal may bias the background estimation from the background-only model, while the background estimation from the background+signal model should remain unbiased. In simulation, this difference is negligible by construction owing to the optimisation of the CR boundaries. When fitting the data with the two models, however, a small difference between the background-only model and the background component of the background+signal model can still exist. This difference is taken as an additional uncertainty. To measure it, the background+signal model from Eq. (4) is fit to the data in the CR and the background component is extrapolated to the SR. After the extrapolation and integration in the SR, the resulting background estimation is compared with the one resulting from the background-only model from Eq. (3). The differences are taken as an uncertainty only in the case of the CI interpretation since it is model-dependent. The σ CRB b uncertainty is smaller than 4% of the nominal background for all SRs of the analysis.
Uncertainties in the signal yield
The expected number of simulated CI signal events in the SR is also affected by theoretical and experimental uncertainties. The signal yield is obtained by integrating the simulated signal in the single-bin SR. This is also performed for all theoretical and experimental systematic variations of the signal. The uncertainty in the signal yield is obtained from the sum in quadrature of the differences between the yields obtained in all variations and the nominal yield. Both the theoretical and experimental components of the signal uncertainty are determined as discussed above for the background in the context of σ ISS b . The theoretical uncertainties, σ Theory s , are presented for reference in Table 3, but are not used in the statistical analysis. The experimental uncertainties of the signal are ≤ 9% for the electron channel and ≤ 22% for the muon channel.
The breakdown of the relative uncertainty in both the background estimate and the expected signal yield is shown in Table 3, sorted by impact. For all cases, the relative uncertainties in the destructive SRs are larger than those in the constructive SRs. This is due to both the smaller size of the SR leading to less background and hence larger relative uncertainty, and the smaller size of the CR leading to a weaker constraint on the background model.
Results
The dilepton invariant-mass distributions for events that pass the full analysis selection are shown in Figure 2. The candidate with the highest reconstructed mass is a dielectron candidate with m ee = 4.06 TeV. The candidate with the highest reconstructed mass in the dimuon channel has an invariant mass of m µµ = 2.75 TeV.
For the statistical analysis, a likelihood function is constructed using a single-bin Poissonian countingexperiment approach. The uncertainties are accounted for as Gaussian constraints taken as nuisance parameters. The compatibility of finding the observed data and the background-only hypothesis is tested by fitting the data with the background model. The p-value of each observation is defined as the probability, given the background-only hypothesis, of observing an excess at least as large as that seen in the data. The Figures (c) and (d) show the region between the SR and CR, but this is not used by the fit. The data points are plotted at the centre of each bin as the number of events divided by the bin width, which is constant in log (m ). The error bars indicate statistical uncertainties only. A few CI benchmark signal shapes are shown, scaled to the data luminosity and superimposed by subtracting the LO DY component and adding the resulting shape to the background shape obtained from the fit. These signals have LL chirality with Λ = 18, 22, and 26 TeV for the constructive case and Λ =16, 20, and 26 TeV for the destructive case. The background-only fit is shown in solid red, with the light red area being its uncertainty. The boundaries of the CR and SR corresponding to the signals used are shown in dotted vertical lines for reference and marked by arrows. The differences between the data and the fit results in units of standard deviations of the statistical uncertainty are shown in the bottom panels. significance is the Gaussian cumulative density function of the p-value. In the absense of an excess, upper limits at 95% confidence level (CL) on the number of signal events in the SR are determined using the profile-likelihood-ratio test statistic [49] with the CL s method [50,51]. These limits are converted to lower limits on the CI scale, Λ. The CL s is computed using 400,000 pseudo-experiments, appropriate for the case where the expected background is small. The statistical uncertainty, due to the observed number of events in data and σ Stat b , has the largest impact on the search sensitivity. The combined likelihood of the e + e − and µ + µ − measurements given the expected background is the product of the likelihood of the individual channel measurements. The signal expectation for both channels is determined by a shared Λ value while the nuisance parameters for each channel remain independent.
The number of events in the SR for the data and the background, and the corresponding significance is given in Table 4. No significant excess is observed. The upper limits on the visible cross-section times branching fraction (σ vis × B) and the number of signal events (N sig ) in different SRs are given in Table 5 and are shown in Figure 3. The expected yields of a few signals, as well as their acceptance times efficiency values in the SR, are also given in Table 5. Figure 4 and Table 6 summarise the lower limits on Λ for the different SRs used in the analysis. The observed limit on Λ ranges from 22.3 TeV to 35.8 TeV. Table 4: The dielectron and dimuon event yields for the data, the expected background and the respective significance in the different SRs used in the analysis. The p-value of each observation is defined as the probability, given the background-only hypothesis, of an observation at least as large as that seen in the data. The significance is the Gaussian cumulative density function of the p-value, and negative significances correspond to deficits. More information is given in the supplemental material 2. This includes information concerning the signal shape and its yields, the fit function parameter values, a comparison of the resulting background with the background from simulation and, finally, the evolution of sensitivity to Λ for different data-taking campaigns, ranging from 5 fb −1 at 7 TeV to the results presented here.
Conclusion
A search for new non-resonant signals in dielectron and dimuon final states with invariant mass larger than ∼ 2 TeV is performed by the ATLAS experiment using the 139 fb −1 of proton-proton collision data collected during Run 2 of the LHC at √ s = 13 TeV. A functional form is fitted to the dilepton low-mass distribution in data and extrapolated to higher masses to model the contribution from background processes. No significant excess is observed above the expected background. Upper limits are set on the number of signal events, as well as lower limits on the CI scale Λ. The acceptance times efficiency values for the corresponding signal shapes are provided. The strongest limits are set on the combined LL constructive model. These observed (expected) limits exclude this model for Λ up to 35.8 (37.6) TeV at 95% CL. The ATLAS Collaboration | 8,401.4 | 2020-06-23T00:00:00.000 | [
"Physics"
] |
Rethinking the Formal Methodology ( II ) : Cognitive Content of Relativity ( On the Centenary of General Relativity )
An attempt to epistemological completion of formal-math theories of relativity is presented. Causal interpretations of SR and GR are suggested. The problem to physical gist of gravity is explained as a contradiction of cognition vs. intuition. Gravity phenomena are represented as unexplored peculiarity of basic particles. The gravity constant is deduced from the known parameters of the electron.
Introduction
We can't solve problems by using the same kind of thinking, we used when we created them.Einstein
Critical Remarks, Objective and Methodology
The significance of gravity and relativity theories as well as the tremendous merits of genius Newton, Einstein and other distinguishable classics is undeniable in today's physics.Meanwhile, we shall emphasize that existing works on examined area both significantly famous ones and not so much known ones, mostly remain as formal-mathematical theories, since any of these does not yet answer such natural question as: what is the physical nature of gravity?
The problem requires explanation of causal link between gravity and material substance that needs to be solved to achieve unambiguous clarity on the issue.
Meanwhile, majority of contemporary theorists consider sufficiently correct quantitative descriptions only as the ultimate goal of the research.Moreover, similar formulation of a problem may sound somewhat unusual-lawless in their view.It seems the deal is not in physicists' competency at all, as they learned to answer the questions "how much is it" before the definition "what it is", within adopted ideology and unspoken standards in modern physics.We urge to draw attention to the important fact in the context: The same experimentally known gravity constant is used in Newton's gravity as well as in Einstein's GR and in many alternate theories by different variations; by the same, the physical essence of gravity remains an unresolved mystery of nature, as it is unanswered yet where from this constant has arisen and why its value is it?
Thus, the significance of existing theories may be evaluated mostly by the quantitative agreement of their results with observations rather than by the cognitive promotion in the subject area, that indicates a pure technical character of the known attempts and acting criteria.The problem is basically important and long disputable, to begin without certain presentation.It demands some detailed examination of a long passed way that has brought to the formation of current ideology, methodology and to a certain stalemate in problematic sections of physics in general.Mentioned intention can hardly be attractive to nowadays' theorists; we realize that it might look as an anachronism, as something long withdrawn from practice because of previous attempts.Meanwhile, it is just inevitable in stated task, since investigated questions demand conceptual-epistemological analysis more than technical.Therefore, certain time and good-enough patience of readers are required for mastering new concepts and language which are actually well forgotten old ones.We start with registration of few important guiding points.As mentioned above: a) The theoretical studies of phenomena are equalized to their quantitative investigations in modern physics, adopted by some historical circumstances (math modeling of reality) The problem was widely discussed long ago among distinguished coryphées of physics, as well as philosophers.The question has mostly related to revealing new-unusual quantitative properties and relations peculiar to the elementary particles of substance (quantum relations).The final key principles and the "correct" methodology were adopted through hard disputes, by the decision of majority despite the unanswered questions.Readers can find some detailed criticism and principles that we will follow in Ref. [1].
We shall emphasize the absence of author's intention to announce existing theories and achieved results on the subject as "something wrong at all".Nevertheless, we have seen: b
) The objective of this work is the cause-cognitive interpretation and completion of the studied subject that is investigated experimentally and mainly quantitatively (mathematically), thanks to deserving researcherspioneers
It must give cognitive "body and blood" to formal math theories, transforming these to conceptually complete, real-physical ones, by author's intention.By comprehensible logic, a critical overview of studied object is necessary for such expectation.The subject of methodology is too large to discuss it fully.We can suggest also Ref. [2] (Russian) on the issue.We have seen a nice book of L. Brillouin as the most valuable on the subject Ref. [3].We bring also his wonderful words on significance of criticism and reexamination of views in science; "Fanatical veneration of any theory is incorrect-they are improving!"Thus, the possible representation of the presented work as a kind of "encroachment" of deserving names and their merits will be deeply unfair.Meanwhile, we just need to agree that any authoritative scientist, even with great merits, may be an ordinary man only, uninsured of human mistakes and misconceptions.The development of science has never gone smooth and straight ahead.Nevertheless, researchers often have been forced to return to the rejected ideas, correcting their mistakes.Then we can assert: The periodical overview of passed way in scientific research must be permissible and necessary; otherwise, we may get a confessional doctrine-instead of realistic science.
Coming to a), we need to explain our approach to the mentioned designation of physical theory in general, to be clear what we are doing next.The matter is we are forced to overpass the adopted hard instructions and recipes on the methodology of physics to get some new opportunities in our investigations.The lacks of accepted paradigm of physical science was discussed and has been criticized by authoritative scientists particularly by Einstein.Therefore, we can be extremely short.The adopted approach is considered as a sufficient condition for complete representation and study of reality, in dominant present ideology.Below unspoken opportunities are supposed by the same.
c) We have the ability of revealing the actual picture of reality through experiments and observations only d) Our abstract-quantitative descriptions completely correspond to reality
Thus, the task of finding "enough-correct" descriptions of "real facts", revealed by experiments, seems as the final objective of a good physical theory.The implementation of further, increasingly complicating experiments and creation of a huge system of quantitative descriptions, covering as many possible facts and phenomena are seen as the desirable task of physical science within conformity of presented paradigm a), c), d).
We briefly have depicted above the essence of the adopted formal-math methodology in nowadays physics.Meanwhile, above-described designation of physical theory a) and convictions c), d) can be evaluated as an expression of trivial desire "to simplify works" that, however, causes serious misconceptions and unsolvable situations in result that we are facing today.Some consecutive analysis should be enough to get convinced in above-said.We will examine one bright historical example somewhat related to gravity problem.
The known historical offer of genius Copernicus, on replacement of the geocentric system to a heliocentric, with consequent huge advance in celestial mechanics directly demonstrates the injustice of above presented perception and nowadays paradigm of physical science at all.Let us remember some details on the issue.The early observations of planets and collected data on their movements did not give opportunity to researchers to see any principle in their intricate paths.Let us assume that our ancestors had been satisfied with observed data and they had registered-systemized these as "real laws of nature", in conformity to c).It is easy to comprehend that they would be forced to use certain tremendous-sophisticated system of quantitative description of planets' movement, well conforming to observations, having no idea of the causal essence of phenomena at all.Thus, their "celestial mechanics" would look as some analog of nowadays quantum theory; there would be well tested, working formulas (instructions, tables, diagrams etc.) and full absence of any causal explanation: why does this group of phenomena go namely so?
Copernicus's incredible merit lies in his thought operations-logical judgments that gave a wonderful opportunity to reveal a universal rule in planets' movement that had been observed before as separately different.He has placed an imaginary observer on the Sun and has defined how planets' movements would seem from there.Then, it became possible "to see" imaginary picture of decent paths of planets and the unique rule of their movements, thanks to this judgments.The discussion is about Kepler's laws somewhat generalizing planets' movement, providing their universal description.Newton's law of universal gravity had next huge advance that was obtained from Kepler's laws by technical-mathematical way only (using differential calculus), in a form of short-compact quantitative relation serving as a basis of classic celestial mechanics.There are many examples in classical physics on significance and inevitable necessity of logical operations.Thus, we can emphasize the important fact: e
) The significant progress in physics was achieved thanks to implementation of logical-judging operations, intermediately between experiments-observations and quantitative investigations-descriptions of phenomena
From this and other similar examples, we can remark next obvious conclusion: f) Mostly, we have no possibility to direct observe the actual values and "the right picture" of reality that we trust and accept as "basic law of nature" As we saw in above example, we have no capability of direct observation of the regular-beautiful paths of planets around the Sun, and we accepted their existence thanks to our thought operations and applicabilityproductivity of created imaginary picture-model.Our next predictions and calculations are based on the created model and adopted principles and can be confirmed with new observations in some favorable cases.However, we can also comprehend the absence of opportunity of new experiments-observations that may confirm our predictions and theories at all, conditioned by different unsolvable technical restrictions mainly.Then we can only be satisfied with the trusted model due to its completion, until new facts force us review our beliefs.It is the normal-natural way of the development of science by its long history.Thus, we can state trivial simplifications in the declared paradigm a), c) and imperfection of adopted methodology in result.
We have used definitely distinguished concept and approach to significance and methodology of physical science with considering the above presented criticism and guiding principles.Reader must overcome natural skepticism and unfavorable heavy suspicions to applied approach created due to long historic circumstances.It is the price to get some clarity on the studied subject that has been natural in many similar cases.Therefore, reader's own decision is required here to judge how much trust the presented work inspires, and how much useful it seems.
The opportunity of interpreting quantum phenomena and microcosm in whole with implementation of imaginary-figurative representations and universal cause-effect laws of nature are presented in Refs.[1] [4].We mark common principles of applied methodology and deep correspondence of basic assumptions.Many realistic thinkers repeatedly have called to return back to natural way of thinking and to cause-effect interpretations vs. adopted formal methodology as Refs.[5] [6].However, official viewpoint on the key principles remains long unshakable despite formal declarations "to involve new ideas".Author's approach to the significance of physical science may be expressed in words of wonderful physicist's R. Lindsay who saw the designation of science "comprehension of the essence of things by thinking" Ref. [7] that does not correspond to dominant ideology.
Mathematics and Epistemology
The necessity to clarify the meaning and significance of used concepts and actions to build our science follow from the above outlined approach.The contemporary physical theories may be characterized as quantitative judgments, corresponding to the experimentally established results, as noted above.Therefore, the examination of assumption (1.1.d)) becomes most important as one of the basic criteria of significance of adopted paradigm of physical theory at all.The history of development and abstract-generalized character of math apparatus are known to us from school education that we formulate as below: We evaluate mathematics as a wonderful human creation, serving as a special language and rational tool, providing important opportunity for description and investigation of kinds of quantitative relations inherent to studied subjects.
By definition, mathematics must work under certain logical control as a "language-tool"; it cannot work by itself and always be useful to us, because of the possibility of its unclaimed applications as with any other tools.It means researcher-operator must well comprehend the meaning and clear target of math actions-operations to get somewhat guaranteed-valuable results.We are forced to mention that modern physicists are often guided by pure formal-mathematical demands only, in abstraction from actual peculiarities of studied real physical objects in dominant practice.They are inclined to trust the strong math rules only due to their standard education, looking at the ordinary logic arguments as some "not enough clear things", therefore, as undesirable and not mandatory!The ignorance of logical control in math operations leaves only the way of quantitative considerations by trivial test-error principle, which increases their works and minimizes the productivity, especially in complicated cases.Meanwhile, a careful examination of the mutual relation of mathematical and logical rules may clarify their common roots and groundlessness of seeming contradictions among these.We can be convinced from history in the incomparable success of mathematics as well as of natural sciences in general, namely, in the period when logical and quantitative considerations were applied with organic combinations.As is known, different kinds of quantitative operations and math functions may be reduced to basic concepts of unit & null (1 & 0) and to application of abstract quantity conservation law (1 = 1 ≠ 0).It reflects the general principle of nature: a real "thing" can be created from another real one only and cannot be transformed to "nothing".Thus, all kinds of math equations and operations are based, and these may be reduced to above-mentioned elementary concepts and simple actions (binary numeral system and modern computing technology may serve as simple evidences to above-said).We will remark for our guidance: a) Mathematics is a specialized section in common study system developed as generalized-abstract, separate discipline by using special-rational symbolism.
The coincidence and correspondence of mathematics with the reality and its workability are conditioned by the implementation of the quantitative conservation laws, reflecting universal peculiarity of phenomena in nature. The quantity conservation laws themselves are based on cause-effect strong logic. Thus, mathematics may be linked to natural science with clear initial definitions the real objects and their properties, expressed by math concepts and symbols
It is easy to realize the above, taking into consideration that the same formulas and equations can be used in different cases, depending what meaning is attributed to symbols.
It follows from mentioned: b) Neither the logical considerations themselves nor pure math methods can serve us independently and sufficiently to complete descriptions of real phenomena. Their clear linked application only can serve as a possible effective analytical way of research
Meantime, above-said is not some discovery in scientific methodology; our ancestors actually had long worked by the same principle.The matter is early thinkers had insufficient math knowledge (as well as experimental capabilities).They somewhat have ignored their significance at all, as natural, being guided by logical judgments mostly.
The unprecedented shift has occurred thanks to the opening of Newton-Leibniz differential calculus.Physicists have decided to review previous methods at all, leaning on pure math methods only, ignoring logical judgments as "something traditional-ineffective", being deeply impressed with the unprecedented success and seeming capabilities of new methodology!The short-term valuable results (the known success of quantum representations) had been perceived as weighty confirmation of reformers' decision, in favor of involved formal-math methodology.Unprecedented problems and confusions however, have arisen with the time because of innovation as Ref. [6].Thus, we can observe from above-said: c) The Ideology and methodology in physics resolutely deflected from one incomplete-ineffective to another extreme, due to historic circumstances 1We have shown above that logical considerations are mostly ignored in present formal methodology.To be more precise, they are actually used spontaneously, in some silent-arbitrary manner, because the attempts to build "pure mathematical physics", without any logical considerations at all, obviously, will be an extreme abstraction that can hardly have any significance!Then simple follows the next unexpected and important demand:
d) Logical considerations in natural science must be used either on the sufficient systematized basis-or be excluded at all
The second way will be obviously speculative and can hardly be useful for someone.Then there remains no choice other than first.Based on above-said we accept: e) The experimental results, logical and quantitative considerations must be adopted as mandatory components in the complete methodology of realistic natural science Thus, one of mentioned three basic tools is ignored in disputable sections of modern physics in fact, due to various historic circumstances.The question has been long discussed by many distinguished coryphées and we tried to present briefly the whole importance of the problem to revise used methodology.The reader, himself, may judge the opportunities and productivity of suggested application by following content.We will mention another important fact on this as well: f
) The combination of logic and quantitative considerations provides a new important tool of research, putting necessary restrictions on each other and mutually controlling both applications
The borders and limitations of applications in study process become clear and appear themselves in natural ways with e) deriving from properties and peculiarities of real objects (which is one of the main problems in present formal-math methodology!).It significantly increases research capabilities and decreases unnecessary mathematisation of problems.The simple examples may demonstrate meaning and rightness to the above-said.
Let's mark apples quantity as A, and number of children B. We use operation A/B and not B/A to distribute apples to children, despite the two operations are equally lawful from formal math's point (as a definition ratio of two numbers).The matter is, here we silently considered that A may be fractional and B never can, i.e., 1) we have applied logical restriction, which frees us from examination of second operation); 2) we get two contrary results ±A when we define the radius of circle with known surface.We choose +A, ignoring −A, because we do not use in practice the circle with minus radius (logical decision); 3) we ignore the sizes of two cities when we speak about their distance.We use such approximations to simplify our work, clearly realizing their restriction and relative significance (idealized, thought operations).These simplest examples demonstrate whole triviality of the supposed opportunity (1.1.d)) and the necessity of initial consideration of physical peculiarities of real objects, parallel with the quantitative operations.It provides the necessary conditions and important instructions to quantitative operations in the researches.Mentioned demands were mostly regarded in classical physics in natural ways (without special declaration), and are mostly ignored in formal methodology due to formed circumstances (also without declaration!).The abstract math concepts become confused with the real-physical ones in the research works that created artificial problems and aroused unsolvable paradoxes in consequence.
The clearly formulated demand to divide physical and math concepts from each other reader can find in Ref. [3] (relating to observation frames, particularly).Einstein has resolutely demanded in his disputes to build physical theories on the conceptual basis, and, the used concepts to connect with the real objects as Ref. [8] that we see in the same context.Mentioned demands however, have met a hard criticism and decisive resistance of majority of theorists and physics has deviated to a present formalism as a result.We shall rely on the ideas and demands of undeniable founders of physics on methodology that we see unfairly rejected.
Outlined remarks and approaches serve us in further examination of study subject.
Physical and Cognitive Significance of SR
Small is the number of people who see with their eyes and think with their minds.Einstein
Causal Interpretation to SR
GR provides certain amendments to Newton's gravity, mostly confirmed experimentally, pointing on its significance and superiority.Meanwhile, logical problems related to unknown physical essence of Newton's gravity were aggravated more with introduction of new unclear categories.These have risen from linking SR to gravity.We shall notice certain improvidence in Einstein's initial approaches with cognitive viewpoint that characterizes the present formal methodology in general and plays a key role in further conclusions.The problem lies in usage of cognitively uncertain concepts to develop new theories.As is known, SR contains some unresolved logical clouds and paradoxes remaining as subjects of hard disputes at present.It is possible to comprehend however, that involvement of unclear categories for solving current problems may complicate them much more by adding an unexplained object to the other dark one (long-term problems with relativity theories evidence it!).Thus, we need to clarify the cognitive meaning of SR before examining GR.We begin with examination of known disputable questions, supposing reader's acquaintance with the subject.
1) Twins' paradox
Travelling brother remains younger in his spacecraft on relation to homebody because he undergoing accelerations that puts certain asymmetry in their conditions; it looks enough-basic to assertion that namely traveler will remain young in relation to homebody, within accepted SR interpretation.However, next simplest objection is possible.We can put symmetry in the experiment by using triplet of brothers for example.One of them can stay in home and two others we will send to travel on contrary directions.Then, it becomes impossible to preferring one of situation for traveling brothers and to decide someway who of them will be old or young?Each of them can calculate by SR principles, using Lorentz transformations (LT), concluding that his brother remains young and not he!Paying attention the viewpoint of homebody also, we fall into deepest confusion as becomes just impossible to find any decision that may be common-acceptable to everybody.It is one of basic criteria to objectivity of science to which SR is not corresponds by its present interpretation as shown above.Similar subjectivism and logic objections are much that pushes many thinkers to reject the significance of SR at all despite some its results are used in engineering level (as ). Described reality demands some not trivial approach, does not sacrificing kind of arguments in a favor to others as it takes place in adopted interpretation of SR.We will look some new examples as well, without quantitative operations, demanding clear answers and giving some hint by the same-where need to look for the causal explanations to logic questions.
2) Change of time and length units with movement
The astronomers had known about velocity of light long-before SR and they well realized that observed picture of any far object corresponds to some of its early state that may be significantly different from the actual one.Then we can comprehend that our brother-twin in the far planet will seem young to us same as he will see us because of limited speed of light.There is nothing mystical here; we can realize that a certain time is needed for light to reach us, which simply explains the phenomenon of observable time difference on distance.Let us imagine someone who travels from our place to our brother.We can realize that in the end of the way his watch will correspond to the brother's watch and his life.Then, it is possible to conclude that his watch will look like "slowed down" during movement to "compensate" seeming difference of time!The seeming correlation and "dependence of time on speed" of traveler becomes clear with above-said; the observable course of traveler's time must slow down more with increased speed of his movement from us.We will also see some distortion in the length of things in movement process, in correlation to speed.We see the two ends of moving meter not at the same moment because of limited speed of light; therefore, its length will look distorted to us, depending on velocity and direction of its movement.We can conclude also that all above judgments are symmetrical and reversible for brothers: particularly, if the traveler moves to us from brother, then his watch will look accelerated to us, and it will seem opposite to brother 2 .We can see nothing against logic in these thought experiments and unusual conclusions as we realize well that the discussion refers to observed values and not to actual ones, which can be different at all.Readers have the right to ask a question here: what do we mean under term of "actual values" in this case?We find the answer from (1.1.f)) in analogy to planets' orbits and Copernicus' wonderful lesson.As we saw, he had used the imaginary observation frame and he took the imaginary picture of orbits as "actual".
Thus, we can accept as "actual" the values and imaginary picture of phenomenon that may be observed (measured) if we will be able to realize instant measurements.It sounds very strange, of course, as we will never be able to see such reality!However, we can create it through our judgments and calculations.We have the same right do it and to trust our conclusions, as we unequivocally believe in the existence of decent orbits of planets in present time.We often use similar actions in classical physics actually, when we equalize friction forces to zero in some cases, for example, well realizing that we cannot do it at all, as it is only an idealized representation (1.2.3.).We just need to take care that we see phenomena not instantly as we silently have accepted (i.e., we cannot take the velocity of light as infinite, same as friction in zero).To get the actual picture of examined phenomena using necessary corrections, we need to remember; the reason of distortions (errors) is the limited speed of light, thanks to which we see "time difference on the distance" and "dependencies of units of time and meter on speed".Then, it becomes simply clear that "the time difference on distance" actually is not depending on speed or on the form of the way but on distance only.We can suppose that the traveler moves by different routes and variable speeds, even exceeding the velocity of light; the ultimate result will be the same because: a) The actual factors causing SR phenomena are the velocity of light and the distance to observed point.Thus, it becomes unimportant how to reach there 3 .
Let us imagine now a researcher who does not realize the reason and physical essence of described phenomena.Then he observes and opens "changes" of time course and length of things in parallel (in correlation) to the speed of movement and he interprets these as the "real laws of nature, revealed by experiments" (1.1.c).Then he declares; "the speed of movement causes the actual changes of physical values"!He makes some quantitative judgments also and discovers certain formulas, corresponding to the results of his observation defining the "actual dependencies of physical values on speed".
Meantime, he faces some logical complications with adopted representations (that we can comprehend already from above content!).Particularly, he cannot answer clearly, what will happen to the traveler if he will move faster than the velocity of light?Then he supposes-declares for such cases especially (to be free from next huge complications): "the velocity of light is maximal in nature that cannot be exceeded any way!"Meanwhile, he has deeply changed the meaning of the basic concept of "speed" by the same.The matter is: b) The concept of "speed" is related to two objects, thus, it can be only relative.Meanwhile, in SR it silently acquires some independent-absolute significance Our researcher faces many similar curious questions and he does not find any explanation: why things must be so contradictory to logic.Then, he decides to "close his eyes" on the logical arguments as "some non mandatory things in physics".Meantime, his theory works successfully despite many logical flaws, since quantitative results mostly correspond to observations that look as weighty evidence to its rightness!Thus, we have suggested above a cause-logical explanation to essence of SR that we continue to examine further.
3) Mechanical speed, light velocity and "space-time" One historical reason to introduce LT and create SR is linked to measurements of the velocity of light in relation to hypothetical environment "ether" that was profoundly different from the expected ones (see Michelson-Morley Experiments.)It was confirmed: c
) The measured velocity of light is invariant, independent from the movement of the observer and the source of light relative to each other
This result contradicts to Galileo's relativity principle (GRP) in fact, as it breaks the rule of speeds' summation on which Newton's mechanics and classical physics are based.We will examine below one extremely simple thought experiment that demonstrates the rightness of mentioned allegation, demanding its clear explanation, not allowing a cover-up of the question with math manipulations.
Let us suppose the experimenter defines the velocity of platform using some standard gun, fixed on it, by measurements of bullet velocity (Figure 1).
The velocity of bullet V is known initially.Its measured value will be m x V V V = + and speed of the platform may be defined as x m V V V = − , according to GRP and summation rule of speeds.We need to notice initially that the same result is possible to confirm by direct measurements of platform speed, without using the gun and summation rule of speeds.Such opportunity and coincidence of two measurements allow us to adopt abovepresented result as correct, corresponding to reality.Thus, we can accept GRP as doubtless, and as a base to other judgments in virtue of it.The creators of SR, also well comprehending its fundamental significance, have announced the correspondence of SR to GRP.Mentioned assertion however, actually remains a verbal declaration only, containing deep internal contradictions that we will show further.
Let's suppose now the observer has used the light source S together with gun.He measures light velocity and got m V c = , according to some experiments and to basic principle of SR.He got 0 x V = if he used GRP and the same summation rule.Then, he decided to follow SR exclusively, as this result does not coincide to reality and it obviously is wrong.The SR formula of speeds summation, corresponding to condition of experiment, is below: + that shows the problem is irresolvable, as x V may have arbitrary values.Moreover, the second experiment does not allow us to tell whether the platform is moving or no! Thus, we can surely mark: The light signal cannot replace a bullet and play the same function of the tool of measurement, independent of our initial convictions or used interpretations.
This fact shows certain qualitative difference of the velocity of light and mechanical movement that actually breaks GRP.Thus, we need an answer to an important question: d) What is the difference of the velocity of light and bullet movement that does not allow their replacement?Meantime, it is only the first part of the problem related to the velocity of light.There is a second important question too that we have formulated below:
e) How the velocity of light becomes invariant, independent of relative movement of the source & observer that directly contradicts to the known rule of summation of speeds?
We must exhaustively answer the questions d) and e) to comprehend the cognitive meaning and physical essence of SR.As we see from (2.1.1)and (2.1.2),SR actually offers below interpretation to the relation (2.1).A new hypothetical participator is actually supposed-introduced in the studied phenomena by linking the "time" to "coordinates" as the "space-time", attributing hypothetical properties to it "to change the physical values with relative speed" as it is necessary for explanation of the observed results.Thus:
f) SR actually considers the inevitable errors of measurements, arising from restricted and constant velocity of light that we are forced to use as a tool of study. Mentioned corrections however, are in fact attributed (verbally) to a cognitively uncertain hypothetical category "space-time" that is represented as an independently existing reality by the virtue of supposition of its own properties "influencing the real physical objects".
We showed (2.1.2) the non-necessity of attributing unexplainable properties to our measuring tools (clocks & meters) "to change their values" with mystical-subjective dependence on relative speed.Thus, the introduction of the concept of "space-time", as a kind of unity "carrying own properties" may be evaluated as a free creativity (inputting a hypothetical reality) that causes explainable logical confusions.Thus: g
) It is important to evaluate the false-fictional gist of "space-time" to remove it from physics, as a directly obstructive factor in problematic subjects related to it
Question concerns to gravity problem in first place, after SR, where "space-time" plays a key role, as well as to physics of elementary particles, where it is involved under the modified name "physical vacuum", with new additional "properties-obligations", necessary for "explanation" of phenomena in this complicated disputable area.
The speculative essence of "space-time", as some kind of independently existing "unobservable reality", may be easily comprehensible if we only agree to take into consideration the obvious-undeniable facts.As we can be convinced, the single natural constant c only quantitatively characterizes full spectrum properties of "space-time" what we actually see in SR formulas 4 .Thus, the whole significance of "space-time" can be reduced to the consideration of the velocity of light in our measurements as it is done in SR in fact; and we already have an exhaustive answer to question: why we must consider light velocity in our formulas (2.1.3.f)).
We bring one additional argument also on the false essence of "space-time".The known combination of three-dimensional coordinates (i.e. the volume) we meant initially under the term of "space".Thus, next realistic question becomes lawful: "about coordinates of what real things are we talking?".The same is right for the concept of "time" too, as we cannot define the "time" (its course, or the interval between some regular events) without using corresponding material objects.These judgments show that: h) The concepts of "space" as well as "time" can be comprehended as the attributes-properties of real material objects; these cannot have physical meaning by themselves-separately, as well as in some of their combination (same as, the concept of "speed" does not have meaning by itself without pointing the objects it relates).We need to mention for justice that Einstein had noticed the meaninglessness of the concept of "space" separately from material objects as in Ref. [9].This judgments and remarks show that the famous innovation of H. Minkowski combination of "space-time" has neither cognitive nor physical significance, if we wish to keep initial meanings of used terms.It will be an obvious nonsense to say; "some kind of combination of properties has its own properties" even from morphological viewpoint.Thus, the concept of "space-time", without mentioning material objects these belong, may have only verbal psychological significance.It creates the psychological impression only to remove the necessity of experimental confirmation of the reality of Lorentz "ether" that was demanded with its definition 5 .The "space-time" has brought a whole group of cognitive mysteries with him as well, on which several generations of thinkers have been working untiringly!The "ether" however, silently continues functioning under new name, because the question: how physical values vary with relative movement needs an answer, same as before.The mentioned fact pushes many researchers to attempt to recover the forgotten "unobservable ether" in modern physics, as the "space-time" plays the same role in its actual interpretation.
4) Problems with dual character of light velocity
As shown above, the introduction of LT and "universal ether" that actually was replaced with the "spacetime" in SR, were conditioned by properties of light as kind of physical reality ("el.magfield" in generalized name) that plays some important role in our measurements and in our world in whole.Our problems with light velocity and its difference from mechanical movement we have divided on two, (2.1.3.f)) that are its restriction & invariance as mentioned in stated questions (2.1.3.b)) and (2.1.3.e)) as well.We will pay attention first on the certain difference of light speed from mechanical.We know that "mechanical speed" relates to a two objects equally and symmetrically; i.e. "mechanical speed" may be defined as the common property of two objects.Meanwhile, light velocity may be defined by different ways: as an individual or, own attribute of el.mag field's exclusively, in first.It is the wave propagation velocity, defined by own parameters of field only: where: λ, T, ν are wavelength, the period and frequency of light accordingly.(2.2) calls also phase velocity.Thus, we can emphasize that V P is the exclusively own wave character of field that deeply different of "mechanical speed" by the same 6 .The velocity of energy transfer by wave group is accepted as the second definition of wave velocity that corresponds to a classical movement of particle; it calls also wave group velocity: where: L, t, are the distance and measured time of wave group's motion, accordingly.We need to emphasize that in second case the etalons of length and time i.e. the used tools of measurements, are independent; these are not defined by own parameters of field as in first case, and these are introduced externally.Mentioned circumstance is important to answer how the invariance problem of light velocity arose.We know that values of both definitions of light velocity coincide for the vacuum whereas wave dispersion is absent.
The same significance of two concepts of light velocity at all is adopted usually for the vacuum, due of equality (2.4).We invite to attention however, that mentioned equality is correct within absence of relative movement between wave source and observer; it becomes impaired in their relative movement.This allegation is easy provable in desire.We just need to consider the obvious fact that Maxwell's equations describe el.mag wave behavior only.Thus, LT for the Maxwell equations have generalized these for the waves only whereas wave source and observer' frame are moving in relation each other.As we see, the application of LT for the Equation (2.2) becomes reduced to transformation: where: ( ) is the universal coefficient of LT.As we can see in representations of ST as Ref. [10] and Ref. [11] γ relates only to Maxwell's equations by fact, because of LT are proven based on the wave Equations (Maxwell-Hertz equations).The matter is λ & T in (2.5) are own parameters of wave field that become changed for the observer due of relative movement in same significance; their relation remains invariant in observer's frame due of it.Similar changes take place of field force-vector parameters also as example: EB , and: Based on (2.6) and definition of rotor in Cartesian frame from Maxwell's first equation in differential form we can write: Or, And we see: (2.7), (2.9) show: light propagation velocity is a natural constant, that is defined by spatial and time changes of force vectors of field; it is exclusively an own peculiarity of the el.mag.field.Then it becomes comprehensible the possibility of expression and detection of this constant in different phenomena connected to el. magnetic field and wave.Thus, the groundlessness of attribution LT to mechanical movement becomes obvious.We must notice that the same explanation to invariance of wave propagation velocity (2.5) actually is contained in Ref. [3], where pointed on "wrongness to take in consideration the changes of wavelength and its period (frequency) together, same time".We see useful to mark also that constant propagation velocity is peculiar not only to el. magnetic waves, and it takes place for mechanical waves too, in particular case, namely when wave source moves in relation to observer.Mentioned asymmetry arises because of participation in phenomenon the environment (that absence for el.magnetic waves).Thus, we can comprehend the lawless of replacement in our experiment the gun with light source (2.1.3.c)) if we take in consideration the wave properties of light (that accepted in most known experiments).
Described explanation is elementary provable in frame of Galilean relativity and classical physics concepts without inputting of hypotheses "changing time course and meters with movement".It is possible to illustrate creation of LT as consequence of confusion and misinterpretation because of we use wave group velocity (2.3) in place of wave propagation velocity of light (2.2).Above-said easy comprehend if use the "particle" properties of light in place of wave parameters and properties.Then we can measure the energy or the impulse of light photons and define velocity of light source, comparing measured values with their initial values in rest condition of the object, as example.We can detect exact coincidence of different experimental results, using light signals and the gun, to define platform's speed in the examined experiment, and not necessity of a hypothetical "space-time", by the same.We use one more example to demonstrate how LT may arise as result of confusion two kinds of light velocity.We need forget initially about "difference of time course", "change meters with movement", as well as artificial problems of "clocks synchronization" etc. that arise with SR in consequence of generalization of invariance principle light velocity in relation to mechanical movement also.Thus, we will remain exclusively in frame of GRP and classical physics in our judgments, by the same.We examine next though experiment to demonstrate above-mentioned opportunity (Figure 2).
We suppose necessary conditions in experiment allowing ignore the signals length and its action time to relation of measured values; it means we can look light signal as a moving point-particle.The fixed time by timer in case of rest condition of the rod, will be: We define the time of reaching the signal to right end (1-2) in moving condition of the rod: 1 Using (2.10), we get: ( ) We define the time to reach signal back left to timer (2-3) by same judgment: 2 , and: ( ) Summary time to passing signal "ahead and back" of the moving rod will be: Thus, the relation of two measurements will be: ( ) It shows that measured average value of light signal velocity on directions "ahead & back", (by using "particle" properties of light) in moving frame to relation of light source, will some less compared to its constant wave propagation velocity by the certain factor: The matter is light signal passes long way on the right (1-2) with low speed (c-V) and short on left (2-3) with high speed (c + V).The average sum of time becomes more with the same, than it has in rest condition of the rod.It corresponds to some slowing of light velocity in moving frame, as (2.15).Let us remember that in experiments prior to ST (Michelson-Morley exp.etc) the light velocity was considered mainly as average of two opposite directions; it are noticed in Ref. [3] as well.Such experiments on direct measurements light velocity (Figure 2) in enough accuracy, as well as on the one direction, actually not implemented (the corresponding links author no succeeded to find).Therefore, these cannot be excluded; we hope it can be implemented on future.Meanwhile, the necessity to explanation of deviation (2.15) from GRP and classical physics has arisen due confusion of mentioned two concepts to light velocity as it shown above.As we know already, problem actually was "resolved" in SR by attributing to our measuring units "time" and "length" the mystical properties "to changing their values with relative speed", in equally, by half for the each!We can to represent (2.15) as below, agreeing with above interpretation: is the known universal factor to "time" and "length" transformation in LT.Thus, the appearance of LT and "space-time" are possible to explain as consequence of misinterpretation, without going out from classical physics and a new hypothesis, if we review our long-term convictions, which is firstly a psychological problem.
We can use the opposite judgments also to show that the whole significance of LT may be reduced to consideration the factors of light velocity and distance in observations of phenomena connected with movement.The existence of certain invariant of "events interval" are exhibited in most of narrations of ST as "undeniable proof" necessity to application the "space-time" (as a kind of unobservable reality in fact): where: s called "events interval" (or, space-time interval); ∆r, ∆t and ∆r I , ∆t I are differences of spatial and time coordinates between events in two frames accordingly, moving relative each to other.The invariance principle is written as below in Cartesian coordinate frame: where: ∆t, ∆t I and ∆x, ∆y, ∆z, ∆x I , ∆y I , ∆z I are time and space coordinates differences in two relatively moving frames.Assuming 0 Applying ( ) ( ) as per LT, from (2.19) we get: (2.20) shows: the light velocity and distance are actual factors that causing difference between results ST and classic laws to the phenomena connected with movement.
Meantime, we saw that the same factors brought to observed distortions of reality in consequence of the movement (2.1.2).Equal results of these factors and interpretations show that the dealt relates to the same phenomena.We need remark that similar explanation to creation of LT and ST as result of misinterpretation of averaging the light velocity on directions "ahead and back" contains in Canadian astrophysicist Paul Marmet' paper also in Ref. [12], next with mentioned remark Ref. [3].The reason and necessity to removing unnecessaryunclear concepts from our lexicon becomes obvious with the same.
5) The "space-time" and non-Euclidean geometry
Other mystic attractive concepts and terms have appeared in different divisions of natural science with creation of SR.We would briefly examine the gist of an important one of these by using previous judgments and conclusions.One of the most known affirmations of SR concerns to the oldest science of geometry.It has been declared in SR somewhat different from "Euclidean" that was long believed as doubtless.Reader can easily comprehend the essence of the question from the previous content.The matter is Euclidean geometry is built on certain axiomatic basis where the static system or, our possibility of instant measurements is supposed unspoken.The issue is that possibility of instant measurement is silently accepted in classical physics as an idealized priory supposition.We can comprehend that it does not change anything in our observations in the static world, i.e. if studied objects do not move in relation to the observer (or, the relative movement insignificantly "slow" compared to the velocity of light).Thus, Newton's mechanics as well as Euclidean geometry are based on the mentioned silent convictions.It becomes comprehensible that their basic principles will seem as "somewhat distorted" in dynamic world, i.e., when the measured objects move with relative speeds compared to the velocity of light.It is easy to comprehend that the observable deviations will depend on two factors.These will be the in-crease with the speed of the movement and decrease with the information transfer speed; thus, the relative difference will look as some function A I /A = f (V/c), where V is the speed of object; c is the light velocity.The above examined though experiments have clarified some aspects of "distortions" of real values pointing on the essence and circumstances of creation of SR.We can comprehend that inevitable deviations of observable results from Euclidean geometry will arise with the movement of studied objects (i.e. in dynamical system) because of the distortion of measured distances and objects' locations relative to actual ones (2.1.2).We already know the causal character of the difference between observable and actual pictures of reality (1.1.f)).Therefore, we can comprehend that it is not necessary to declare Euclidean geometry as a "conceptually wrong" science that needs to be replaced by some other kind.Meanwhile, we must just consider the limited speed of measurements (observations) that causes the mentioned differences between the observable and actual pictures of subject phenomena.Then we get new description rules and geometric laws where the inevitable errors of observations are considered.We can call those as "pseudo Euclidean" or "Lorentz geometry" if desired, comprehending however, that it is actually the same Euclidean geometry, where the errors of observations are considered.The "new geometry" gives us opportunity to make the calculations and get results that coincide with observations.However, we lose another important capability with the same.The matter is, we cannot use the description by "new geometry" for the cause-effect (or logical) investigations of phenomena because it contains the errors of observations What do we need to do in such situation?Copernicus already gave the exhaustive answer in time that may be applied to this case also (1.1.f)).We must just clearly divide the values, descriptions and pictures into observable and actual ones.Then we can understand where and how to use each of these correctly as well as how to pass from one kind of values to the others.Namely, if we need to investigate the cause-effect side of phenomenon we must recover the actual picture from its observed one, by taking into account the errors of observations.These arise because of limited velocity of light, in context of the studied problems.We can recover the picture of the phenomenon that we can see if the velocity of light will be infinite.The picture that we get in such a way will correspond to its description in the idealized Euclidean geometry that can serve us in the cause-effect investigation of phenomenon.However, we need to "go back" again to "Lorentz geometry", where limited velocity of light is considered, to get the opportunity to compare our conclusions to experimentally-observables.Thus, the task and operations of such transformations from one kind of geometry into the other are principally similar to that used to study the movement of planets, using the geocentric-heliocentric-geocentric systems transitions, the significance of which is clear to us and does not call for any questions.We just need to realize the flaws of our observation system in one case, and the imperfection of our measuring tool in other, which we must consider in our actions and judgments to get correct conclusions.The matter is we have put the questions in different ways in formal and in realistic methodology.In the first case, we wish to have the description of the phenomena, which coincides with our observations, as it seems to us, without thinking of its cause-effect side.However, in the second case, we wish to penetrate into cause-effect essence of the studied phenomenon.Then we must do some additional operations with our results of observations to "filter" them from inevitable errors connected to imperfections of real systems and measurements.The imaginary ("clean") picture of the phenomenon in the idealized system of observation only can to serve us for its causal investigation and correct conclusions.
6) On the significance of SR We have shown above that the cognitive problems of SR were arisen due of trivial confusions of used concepts as well as with the arbitrary interpretations.We hope this explanation may be easily perceivable, despite the whole painfulness and the huge psychological problems related to suggested recognition.
i) We see the most negative role of SR in introduction of the uncertain concept of "space-time" that has prevented the development of subject divisions in physics
Meanwhile, we see inexcusable to announce SR as some "reactionary and totally harmful falsification that needs to be excluded from natural science at all" as demanded by most critics and hard opponents of Einstein's theory.We can emphasize next, with all criticism and marked flaws, some of the new visions and non-traditional innovations of SR that give undeniable push to resolving certain huge problems as:
j) SR provides quantitative descriptions of certain phenomena that mostly coincide with observable results; it may be used at applied-engineering level. E = mc 2 fundamental relation is one of the undeniable huge shifts in natural science
We also remark next unprecedented innovation of SR from methodological viewpoint:
k) The description of phenomena in different frames of observation with comparison of their results gave principally new opportunity to reveling unknown relations of nature
The opening of mass-energy communication was possible to reveal thanks to studying the same phenomenon in two systems of observations, as for example.Thus: We need to perceive SR within its actual significance: as a way of description of phenomena where our real capabilities and admitted inevitable errors of measurements-observations are considered, next to cognitively irrelevant interpretations that need to be replaced by cause-realistic ones.
Physical Essence of Gravity
I believe The Lord has decided what we need to understand and what not, but allowed us to try! Author
Acquaintance to Stated Problem
We ubiquitously see free falling of things and feel their weight, perceiving these quite ordinary, not having any idea and easily neglecting the causal essence of the oldest mystery of nature, called "gravitation".Such preamble may sound outrageous and unexpected for many, on the background of frequently launched rocket-satellites at present with applied or research purposes.We also periodically have learnt about next confirmations of this or that predictions of Einstein's famous theories and of other remarkable achievements, closely related to examined subject demand a serious theoretical base and sophisticated calculations, inspiring the opposite impressions.Meanwhile, the complete darkness of physical gist of gravity is today's reality in natural science and our reader must take this allegation seriously, resolution of which is one of the main tasks of suggested work.The stated problem somewhat dropped out of attention at present, in comparison to researchers' early attempts devoted to open the physical essence of the gravity.Present efforts and huge means however, are mostly directed to study of the quantitative side of the gravity phenomena in fact, due to certain historical circumstances, as noted (1.1.):
a) We strive to define the causal basis of the examined phenomenon initially, from which its quantitative peculiarities can be derived naturally
The Great Newton actually did not tell anything on the physical nature of gravity in his time, proving just that we get the right description of the observable movement of material bodies and the known celestial mechanics supposing far-acting forces in Inverse Square of their distance (1.1).Newton's gravity does not give answers to natural questions: what is the nature of that force?Through what environment and how it passes there?Meantime, the mystical force, instantly acting on unlimited distance, looks unnatural and it was perceived by theorists skeptically from the beginning.The concept of "gravity field" was introduced in physics with Newton's gravity, as a "transmitter" of the hypothetical far-acting forces in fact.It has been presented uncritically within close analogy to Coulomb's electrostatic field in most textbooks, using similar terminology (as the "gravitational potential").The "gravity field" has been looked as a kind of physical reality, as "it having its own peculiarities", that has only hindered the actual cognitive problem for a long time.
We can look at some of the new theories with involvement of Mach's Principle as contemporary modifications of far-action and field-based theories of gravity in which the instant action has changed with finite speed of gravity influence propagation equal to light velocity as Ref. [13].Mach's Principle however, must not be acceptable for us since it also obviously supposes the existence of a new kind of physical reality (as influence transmitter) without any experimental evidences of its reality (see point c) below).We need to emphasize however, the mentioned approach gives quantitative results equivalent to Einstein's GR.By the same, it may serve to us as an additional testimony of the causal essence of gravity in our further attempts.
b) Le Sage's theory of gravity was one of the conceptually formulated explanations of the essence of gravity, from the contemporaries of Newton.It was based on the existence of kind of special hypothetical particles that move in all possible directions in space with much more speed than light velocity, arousing the screened (or shielded) effects between material bodies, being partially absorbed in these.Le Sage's particles however, did not stick because of different serious counter arguments, one of that pointed by Poincare Ref. [14].
c) We will briefly examine here the contemporary modification of "special particles" theory of gravitation which is related to introduction of mass-less gravitons as the mediators of force transmission at distance in quantum field theories.Many modern works have been developed with the gravitons now.Meanwhile, direct detection of single gravitons experimentally seems practically an irresolvable task because their energetic insufficiency.An indirect confirmation of their existence by detection of gravity waves (as their coherent groups) seems realizable; practical works on this direction are underway at present (see LIGO, VIRGO, LISA experi-ments etc.) We will not examine the technical base and disputes on this matter referring to existing large literature on the subject as example Refs.[15] [16].However, we shall pay serious attention to the fact that experimenters have been looking for gravity waves for significantly long time.The used techniques are much improved (with the costs) starting with earth-based antenna-detectors and passing to cosmic interferometers having incomparable capabilities at present.The gravity waves however, remain undetected yet, despite the achieved fantastic sensitivities of detectors.One alarming factor here is that the different significances of gravity signals are accepted within time.The suspicious conclusion simply follows then without touching the technical details of the problem at all.The theorists are doing their calculations issuing from certain assumptions that give certain initial data.The experiments, however, do not confirm their predictions in fact.Then, they change the initial data and their basic supposition too, to get other result than the earlier adopted one (in this case much smaller energy of gravity waves).It shows that apparently, theorists are not even working by test-error principle (1.2), and they just strive to adjust their calculations to the results of experiments i.e. without any definitive base concept.Mentioned circumstances do not correspond to initial criterions of objectivity and methodology of realistic science that have been declared.Therefore, we can appreciate them as unreasonable expenditure of efforts and means.It seems appropriate to recall here the practice of introduction of hypothetical kinds of realities to explain the incomprehensible phenomena, largely used by early thinkers in complicated cases; these mostly brought nothing but irresolvable mysteries.The known histories with the "phlogiston" and with kinds of "ethers" may serve as examples.
The Great Newton said: "Hypotheses non fingo" (I contrived no hypotheses).We can interpret this famous expression as a transparent commandment-do not harm natural science by own compositions!It seems as pertinent to refer to the similar opinion on harmfulness of "unnecessary essences" (see Occam's razor) as in Ref. [17].Mentioned principles express the reasonability and mandatory the preference in natural science of confirmed facts relative to arbitrary suppositions, by authors' interpretation.Thus, above-examined critical remarks and accepted methodological criterions gives to us full rights to reconcile with lack of gravity waves and to think; what would that mean?
d) The absence of gravitational waves exacerbates problems with gravity much more, at the first glance.From other side however, it may give to us next valuable instruction:
We search the causal explanation of gravity problem on completely wrong direction since recognition of absence gravity waves demands basic changes of the considered versions and representations on the physical essence of gravity in whole.
The existence of "gravitons" and "gravitational waves" accepted to represent as the derivatives from GR. Meantime, the famous theory of Einstein-GR remains yet without any of unambiguous causal explanation as noted; it gives large possibility to interpret this or that its quantitative results in arbitrary manner, by theorists' personal propensity (we see it actually with presence several different theories in similar quantitative significance!).
The tremendous cognitive revolution has been predicted by few authoritative specialists on this area, needed to open physical essence of gravity phenomena as Ref. [18].
The GR (and equal theories) are adopted by specialists as the adequate quantitative description of gravity phenomena at present, most results of which are confirmed experimentally, excluding the "gravity wave" that we will discuss.
Authors share the majority's opinion; Einstein's theory gives sufficiently correct quantitative descriptions to many observable effects, related to gravity phenomena.We emphasize however, that GR remains as completely non-trivial from the causal-cognitive viewpoint, the illumination of which is one of the main tasks in the work.
e) The statement of question is possible to present as the following allegation:
GR (and equal theories) are satisfactorily correct theories in quantitative meaning. Therefore, these must have enough informative content to open the causal essence of examined phenomena as well, in virtue of their correct quantitative relations.
Above-said is right to assert in relation to Newton's gravity as well, with consideration of it as the approximation of GR.The mentioned opportunity and advantage follow from the adopted methodology, with the demand of parallel usage of mathematics and logical considerations in research process (1.2.a)).It allows translation and passage from one kind of language and description to another, upon necessity.Thus, we just need to define the correct physical meaning of used mathematical symbols and actions, to pass to the descriptive language and causal interpretation of gravitation from the correct quantitative relations that are already known to us.
The Causal Side of Gravity
We shall briefly examine some interpretations and terminology adopted in GR and other gravity theories from cognitive viewpoint that may hint on the essence of the problem.We shall remark first that Newton's "instantly far-acting" force was removed in GR that we see as an important advance in cognitive meaning: as a rejection of arbitrary hypothesis!a) The gravity phenomena are interpreted in GR as non-trivial consequence of curvature of the "spacetime", as presented in most literature.It creates the following impression and causal picture of the gravity problem, if guided by pure formal logics: A special kind of reality, "space-time" exists that becomes "curved" in surroundings of a central massive body under its influence.The "curvature" of "space-time" acts on test material bodies forcing them move to central body by acceleration (free falling), or be pressed on its surface (weight force) after their collision.
Meantime, we can just state the absence of any experimental results for today that directly evidence the existence of kinds of physical realities corresponding to demanded peculiarities of "space-time".
It is easy to comprehend also that the same facts could serve for proving the reality of Lorentz' ether, Le Sage's particles, as well as the existence of gravitons etc, if these could have been observed in its time.The next formal interpretation also adopted that sounds as: The pseudo Euclidean Lorentz geometry (2.1.
5) turns into Riemannian under the influence of material substance. Spatial geodesic lines are peculiar to it and material bodies move along these free from force influence, in the form of orbital movement.
Then some impressions and corresponding interpretation arise, such as "the Riemannian geometry causes gravity phenomena at all"!We hope it will be easy for the reader to comprehend the absence of any physical meaning in this formulation from previous content.We just need to state here that "geometry" is a kind of science, to comprehend above-said.Then it becomes clear that "geometry" (as well as any science) is a way of description (tool, language, system, i.e. a human's abstract creation) that itself is unable "to influence" the physical reality!We can see here a simple confusion of abstract math and real physical concepts [3], creating a nonsense, peculiar to formal methodology (1.1, 2.1).Thus, we can state the actual absence of any third real physical participant in the gravity phenomena, being observed between two objects, independent from interpretations and used terms (as "special particles", "physical fields", "ether", non-detectable "space-time", "curved geometry" etc.) We see that non-ordinary terminology and concepts of GR are conditional names only to mark some math objects and actions.Thus, it will be meaningless to use these for causal description, because of the absence of their initial physical meaning at all.
b) The actual significance of "space-time" in SR is reduced to consideration of light velocity factor in our measurements (observations) in a form of universal correcting coefficient γ (2.16) as it is shown in the previous chapter.We show above, the movement of studied objects and limited light velocity has caused the difference between Newton's physics and SR.We can get convinced that all of the confirmed gravity effects of GR, distinguishing it from Newton's gravity strive to zero if we accept c → ∞, which means GR turns into Newton's gravity, same as SR turns into Newton's physics.Above-said, however, transparently instructs that GR effects are the consequences of a certain dynamic process, from the logical viewpoint.It just means that gravity effects are conditioned by the movement (as SR effects were).The argumentation of this conclusion is obvious; if we deal with static world and unmoving objects, we will get the same results of measurements independent from speed of our measurements, i.e., the light velocity should not be expressed in experimental results and in our formulas!(It means our geometry always will seem Euclidian).We shall compare some known expressions of SR and GR to show the rightness of this conclusion.The invariance of elementary interval of "space-time" is written in SR as: where, the relations for elementary spatial and time intervals in two frames are the same: We mean v in (3.2) as the free variable by its definition and all of SR effects were simple consequences of movement as shown in previous chapter.It is easy to observe that GR consequences and effects may be represented as similar functions of v/c relation.The linear element for spherical-symmetric Schwarzschild metric, for example, presented as: ( ) where: 1 s r r e ν − ≈ characterized the "space-time curvature", r is the distance of point from centre of material body M, and G is the gravity constant.The The physical meaning and significance of speed ν g will be examined next.We shall mark only that it is not free variable in GR as it is in SR, and ν g defined by parameters of material substance as (3.4).Thus, the identical structures and values of e ν and γ 2 factors evidenced: c) GR effects and the term "curvature of space-time" (or "gravity field" etc.
) must be comprehended in the same meaning as the effects of SR; i.e. these are observable distortions of reality caused by the objects' movement and by limited light velocity
Thus, this conclusion is justified from both sides: from ordinary logical viewpoint as well as from purely formal consideration, in virtue of the same physical values and their same combinations cannot be interpreted in different ways, in any scientific methodology.The possibility of similar transformation (3.2) and (3.4) as function of v/c can be observed in all confirmed effects of GR, without any exception, as for the displacement of planets' orbits, gravitational change of frequency, frame dragging effects etc., that also confirm the above conclusion (examples are shown next).Thus, we come to a clearly formulated next conclusion: -The gravity phenomena are consequences of an unknown to us movement.The known Einstein's Equation (3.5) in GR are perceived and interpreted by theorists as the "field's equations".The structure and the physical units 7 of components in (3.5) simply show, however, GR equations are related to a motion and nor to kind of physical reality as it declared.
Then the "gravity field", "curved space-time", different "special particles" etc. become groundless hypotheses and arbitrary interpretations in fact.This conclusion however, immediately collides into incredible problems.The matter is we do not see any kind of movement that may directly confirm it.We observe, for example, weight of things in their obvious unmovable condition and we cannot imagine on what kind of movement we can talk here.We shall however put aside different "obvious" questions to continue ubiquitously examining our conclusion unequivocally deriving from examined arguments.
d) The "local equivalence of gravity with inertia" (Einstein's equivalence principle) (EEP) is adopted as the other most important basic principle in creation of GR
We shall firstly mark the mysterious character of the adopted allegation from logical viewpoint, that we have the right to discuss, because the used terms "gravitation" and "inertia" are exclusively concrete physical concepts demanding clear definitions of their meaning.Ordinary reasoning tells us that real physical objects can be individually independent things or, these may be the same thing with different names only.It will be an obvious nonsense to say something as; "the objects A and B are individual at all but may be the same things within some conditions!" 8 .We know such considerations in mathematics, for example, by accepting the average value of some numbers as equal to actual, or by adopting the curve element as "straight" in differential calculus.The matter is such approximations have meaning if compared objects are of the same kind, otherwise we will fall into obvious nonsense (as if comparing "mass" and "distance").This demand is preserved in above examples, as compared concepts are both of the same kind ("numbers" or, "lines").The concept of "inertia" we can define as a phenomenon only, arising in consequence of accelerated movement, thus as a dynamic process.Then we can conclude: -We are obliged to consider the phenomenon of "gravity" as consequences of accelerated movement, guided by the demand of uniformity of compared concepts.
Thus, we got one more independent instruction of the accelerated character of unknown movement (3.2.b), c)) that causes the "gravity" phenomena.
We shall present now the circumstances that have pushed for the acceptance of the mentioned strange allegation in GR.The phenomenon, caused by accelerated movement, well known to us from Newton's mechanics.The most important certain peculiarity of the inertial mass of material objects is revealed in the phenomenon, characterizing how the body resists to external force and acceleration.It may be experimentally defined from Newton's second law i F m a = as: where: F is the acting external force; a = dv/dt is the acceleration of movement.
The number of experiments, starting from Galilee and further much more exact ones, has shown that gravity forces acting on the material bodies are defined by another common peculiarity, independent from the kinds of tested materials.The gravity force that acts on the test body is defined by Newton's law of universal gravity as below: The m g in (3.7) is called the gravitational mass that characterizes corresponding peculiarity of test body; namely, it shows how much weight force the body generates under the gravitational influence, the physical nature of which is still unknown to us.Thus, the above-mentioned experiments show unequivocally that these two different experiments give exact equal values for the mentioned two kinds of characteristics of the substance: The equivalence (3.8) brightly expresses all the mystery of gravity phenomena since the natural question arises there why these two kinds of characteristics must be the same that remains unanswered, despite the big number of theories, uncountable written pages and long disputes.Moreover, the reader must know that many other kinds of experiments with the gravity and with the accelerated movement give the same results: -The existing large group of facts shows that the consequences of "gravity" influence are indistinguishable from those that arise in consequence of accelerated movement.
It concerns the above-mentioned phenomena: 1) force influence and movement (i.e.weight and free falling); as well as 2) geometrical changes of light's trajectory; 3) the gravitational change of light frequency; 4) the gravitational delay of time and other effects.Mentioned equivalence of the gravity and inertia are mainly confirmed by number of experiments that reader can find in the literature, as for example Ref. [19].The possibility of representation of the known effects of GR as the consequence of accelerated movement is also shown.Einstein has taken into attention the mentioned experimentally established fact in creation of his gravity theory as a key principle that has been formulated in EEP.We invite reader's attention to the next important remark once more.The matter is:
-Einstein did not explain any way the similarity of the two kinds of phenomena (gravity & inertia) that were adopted as different subjects with initial definitions; and he only stated the experimentally revealed facts, considering these in the base of GR.
Moreover, the genius thinker never hides that he does not comprehend the causal essence of gravity phenomena; he says, for example, "If I could only understand what goes on in a falling lift!" Ref. [8].Thus, from previous content, we shall evaluate Einstein's relativity theories as they actually are, i.e., as pure formal-mathematical ways that provide important results, mostly corresponding to observations.Therefore, the attempts to present our conclusions as somewhat opposing Einstein's relativity theories will be obviously inappropriate and groundless, because of actual absence of any causal interpretations there at all, by definition (1.1 b)).e) We shall examine now the important question: where from arises the restriction of "locality" in EEP? Named restriction only hindered direct identification of two concepts (inertia and gravity), which may signifi-cantly change created situation.The logical vague character of EEP has pushed some theorists to reject it at all, accepting inertia and gravity as separate phenomena having nothing common with each other.In fact, the gravity phenomena are attributed completely to peculiarities of "space-time" in RTG as Ref. [20] ignoring mentioned principle for example.Meanwhile, the attempts on opposite direction have taken place in other works, by presenting "inertia" as a consequence of "curved space-time" also, to explain the similarity of the two phenomena.The presentation of "space-time" as "owner of peculiarities" however, simply transforms it to a kind of hypothetical physical reality, on the rank of "non-provable ethers", the wrongness of which was shown in previous chapter.Moreover, by this approach, a large group of equal results of mentioned experiments (3.2.d)) seems as a number of unexplainable exact coincidences.Thus, from above argumentations we see the inevitable necessity of a clear definition:
Is the deep similarity of gravity and inertia phenomena a pure coincidence or are they identical by their physical essence with different names?
The first choice seems just unbelievable with elementary logic, by taking into consideration the equivalence of inertial and gravity masses only, leaving aside even the large group of other kinds of coincidences.The conclusions on equivalence of gravity and inertial masses of Galilee and Newton had been confirmed by Eötvös with impressing accuracy (10 −8 ) about hundred years ago, later with much more exactness (10 −11 ) Ref. [21], and known last results were achieved in 1999 (10 −14 ) Ref. [22].However, new projects to test the equality of gravitational and inertial masses (weak equality) with unprecedented accuracy (10 −18 ) are suggested at present as Ref [23].Meanwhile, the above-described reality clearly shows that the ubiquitous similarity of gravity and inertia phenomena are accepted by researchers as a statement, as they examined the question "how similar they are", and not "why they are similar".Thus, from cognitive viewpoint we can state: The restriction of "locality" in the equivalence principle has banned direct identification of the concepts of "gravity" with "inertia", which caused further huge cognitive complications.
However, scrupulous examination shows actual absence of any quantitative exposition of mentioned restriction of "locality" in GR.Reader himself can get convinced that nothing changes in GR if we replace the "local equivalence" by direct identification of "gravity" with "inertia".We see in GR m i = m g (3.7) adopted without any conditions or criterions; it shows that the concepts of "gravity" and "inertia" are actually indistinguishable in quantitative meaning.Thus, pure verbal-psychological character of the "locality" restriction becomes obvious from above-said.The mentioned fact is obvious and can be checked up if desired.Some critics of GR have also observed that concepts of gravity and inertia are quantitatively indistinguishable as Refs.[20], [24].Moreover, academician Fock had pointed in his book on EEP; "The law of equivalence of inertial and weight masses have a general and nor a local character" in Ref. [25].We need to emphasize that a number of experiments confirm the exact equality of gravity and inertia but not the "locality".The fictional essence of the restriction of "locality" gets obvious; we can assert: Experimental results as well as quantitative expressions, confirming restriction of "locality" in equivalence of gravity with inertia are absent in GR.Thus, the allegation of "locality" may be evaluated as a verbal declaration, adopted in virtue of intuition.
We briefly examine beliefs and prejudices induced to adopt the restriction of locality in EEP that hindered the acceptance of gravity and inertia as the same thing when facts and reasons to do it are many.We know that the results of experiments in the accelerated spacecraft and in the unmovable lab in gravity field are the same; by free falling of test bodies, by force reaction (weight), by deflection of light path etc. Mentioned similarity of results does not allow the inner observer to define whether his closed lab is in accelerated movement condition or it is in corresponding gravity field?The equivalence principle was adopted due to similar results.However, the solution of the problem and detection of the difference between gravity field and accelerated movement seems to be possible if we use an "enough big" lab (or enough exact measuring tools).The trajectories of falling bodies for example, directed to the centre of material body, being the source of gravity field that are not parallel (Figure 3(a)).Meanwhile, these are supposed to be parallel in the accelerated spacecraft (Figure 3(b)).
The above-described conclusion seems sufficient to put the restriction of locality in equivalence to gravity with the accelerated movement.This conviction however, is completely based on supposition because the difference between gravitational and inertial phenomena is not yet confirmed experimentally anyway.This restriction also does not have any quantitative expression in GR in fact, as noticed above.Then we mark: The restriction of "locality" in GR plays purely declarative-psychological role, introduced because we do not see directly the corresponding movement.
f) The identity of the concepts of gravity and inertia: universal expansion of substance
The concepts of "gravity" and "inertia" are actually used in the same quantitative significance in GR, examined above.This conclusion opens clear indication of the physical nature of gravity.Moreover, it is easy to see that the problem of "locality" by itself goes out of agenda if we adopt the "gravity" as the "inertia", sacrificing our intuition.Thus, the whole group of phenomena inside our terrestrial lab that we call "gravitational", can be looked as consequence of accelerated movement.It follows that our Earth, for example, continuously expands pushing the things on its surface in radial directions with acceleration, and the bodies resist to acceleration with inertial forces according to second law of Newton.The free falling of different kinds of test bodies with identical acceleration becomes simply explainable; the things, being free from their supports actually remain in their former places; the surface of Earth reaches them simultaneously.
However, we will be unable to detect the described expansion visually due to its universal character for the material substance at all.The cabin of accelerated spacecraft as well as our etalon meter will expand proportionally with all other material objects for the same reason.The trajectories of free falling test bodies will be indistinguishable from those, which have been in the terrestrial lab as illustrated in the graphic, i.e. they become not parallel (Figure 4).Such explanation of gist of gravity strongly contradicts to human intuition due to our daily perception of material world that directly hinders even its detailed study.
Many "obvious" objections also immediately arise there, for example such as: how can orbital movement and the celestial mechanics be explained, replacing universal attraction by expansion of material substance?The problem, however, is not new in the history of science from cognitive viewpoint.The humanity, for example, has been forced to agree with the "rotating" Earth and its orbital movement, with incredible velocity that was "the most unmovable thing" for us.We cannot see this, and we have adopted it today as out of doubt-after paying a proper price!Then it is possible to comprehend that we are in the same situation; we need to pay the next huge price, going versus our natural intuition to solve the mystery of gravity.Different kinds of gravitational phenomena then become possible to interpret on the comprehensible causal base, without exceptions, (some examples follow).
The concept of proportionally expanding material universe gives us important evidences of solution of many other fundamental problems of physics that will also be discussed further.Firstly, however, we shall draw the attention to below historical comparison: we have intuitively formed many "doubtless" convictions on surrounding us material world due to our direct perceptions.We were initially convinced in the: 1) Absoluteness of rest and movement; 2) absoluteness of directions; 3) opportunity of absolute (instantly) observation-measurement; 4) absolute invariable sizes.
We learnt the history and dramatic events that forced us to remove the first two points from our minds from school education.We tried to show in the previous chapter that logical problems with SR reduce the necessity of releasing the third point from this list of false convictions.The fourth point however, remains yet strongly unshakable in our mind that forces us to refer to inappropriate creations to save our conviction, arisen because of natural intuition.g) Other evidences on expanding world 1) The causal interpretation of physical nature of gravity is absent in Newton's gravity theory as mentioned above (3.1.a)).Meantime, the actual identification of inertial and gravitational phenomena is silently used there in fact, since Great Newton does not put any difference between the two kinds of masses in his theory and quantitative considerations at all.It simply brings to the same conclusion in favor of expanding world, as it was narrated above.
2) We can find the valuable evidence of universal expansion in Ref. [13].This theory provides correct results equal to GR (3.1 a)) and the conceptual explanation of gravity is based there on the variability of particles' mass.The brief content of theory is as below: All kinds of physical values are possible to represent by combinations of h, c natural constants and with a single basic value only, serving as free parameter, having measures L, T, M, etc.This possibility derives from well-known quantum relations: where, m mass of particle; h Planck's constant; λ c , ν c , τ c are Compton's wavelength, the frequency and the wave period, accordingly.We need to emphasize here that the free parameter may be variable, in virtue of its uniqueness.The laws of nature and observed phenomena in general will remain the same with this supposition (it means we will be unable to see any changes in our world).This conclusion was used in mentioned work by accepting the mass as a variable.However, the relations (3.9) unequivocally say that λ c (ν c , τ c) also must change with variable mass of the particle!Thus, we are just obligated to adopt continuous expansion of Compton's wavelength and decrease of its frequency, if we accept continuous reduction of the masses of particles.It means our world is in dynamically variable condition, i.e. our meters and our clocks etc are permanently in change together with us and with all of material things, in such a way that leaves no opportunity for us to perceive our real situation directly.We observe reaction of forces between contacting material bodies (weights) and we see their "free falling" each to other that remains completely unexplainable to us due to our intuitive convictions!The difference of two theories (GR and [13]) relates to their verbal interpretations only having no actual significance in the results; thus, with the same, the described conclusion on expanding world relates to GR as well.
3) Hubble's expansion of universe (Hubble's flow), that is now accepted as a doubtless fact by dominant majority of experts, may serve as direct evidence to expanding world.
It was established by observations that far away galaxies recede from us with some speed proportional to their distance.This is characterized by Hubble's law; V H ≈ H 0 D, where H 0 ≈ 75(km/s)/Mps (Mps ≈ 3.09 * 10 22 M), and D is the distance to observed object.We need to bring only one important remark on this matter.The dealt is a strange situation created by this wonderful opening, arousing continued disputes among theorists.The expansion of universe actually follows from GR (it is simple to understand from previous content).It was theoretically shown particularly by de Sitter, A. Friedman, G. Lemaitre.Einstein added then the special constant Λ (cosmological constant) in his Equations (3.5) for the sake of "protecting the static condition of universe".However, its necessity disappeared with the opening of Hubble's expansion.Thus, the following picture has been created: The "gravity influence" pushes universe into compression, vs. expansion of Hubble.The intriguing question immediately begs here: which of these opposing factors will prevail in the fate of universe?Then the experimenters have begun the measurements of density of substance in universe for necessary correction of Hubble's constant to solve the arisen problem.The statement of the question is the following: is it the less or exceeding average density of substance than the critical ratio in universe that will define how the expansion will go in future?Will it continue forever or it will stop and change to compression etc.The surprising fact however, has revealed; the corrections show that relation f(H 0 , ρ) → f(H 0 , ρ cr ), which means the factor of expansion has appeared too close to gravity compression factor that does not allow yet to answer clearly what will happen to universe in future.The solution of problem is possible by analyzing the above-mentioned facts.
The universal expansion follows from GR i.e. from quantitative description of gravity. It just means that both of these factors are different expressions of the same phenomena; their equal quantitative significance becomes then simply explainable.
Moreover, the accelerated character of expansion of universe also has been confirmed by observations as in Refs.[26] [27].Thus, we have a complete opportunity to replace gravity by universal expansion of substance in virtue of EEP and previous remarks, and the absence of visual perception of expansion only prevents us to adopt it.
We need to examine one wonderful question arisen with Hubble's expansion, to evaluate the actual significance of our visual perceptions.The matter is there is no reasonable answer to the question; why does universe expand within global scale and is it unchangeable in short-local scale?The question is lawful since the same common laws of nature determine behavior of material objects independent from scales 9 .Then we can simply state, judging from the circumstances of the problem.The image of expanded universe was created using different methods of evaluations.The expansion of universe in large scales was accepted in virtue of Doppler shift of light frequency, meanwhile we judge about unchangeable sizes of our planetary system as well as galaxy in virtue of direct visual observations.We have no technical opportunities to visual observe of geometrical changes in large scale, for the faraway cosmic objects, and we adopt their motion in virtue of light's frequency change only.The frequency changes, however, are peculiar to the local scale cosmic systems also that are visually seen unchangeable (red shift of Sunlight, reaching to Earth, for example).The deepest subjectivity of our methodology becomes obvious from above-said.We observe the frequency change on distance as a common peculiarity of universe, in fact.Thus, we adopt the frequency change as evidence of movement and expansion in large scales, where we are unable to observe visually.However, we explain the same results of observations as an unexplainable to us "gravity influence", in scales suitable for our visual observations.Thus, we are facing the inevitable choice: -We must adopt the observed cosmic expansion as common-universal, attributing it to material substance also, sacrificing our intuition, or: -We must accept a large group of known facts as a chain of incredible coincidences.
The second one was accepted, in fact, due to huge pressure of human intuition.We prefer the first one however, considering the known role of mentioned factor in science history in general and the harmony of world perception that opens up with this.Then this conception becomes well confirmed with dominant conviction of creation of Universe (Big Bang Theory) and it provides additional evidence on the issue.It simply says that all cosmic objects, which we now see in gravitational balanced condition, were created from an insignificant "small" space and from one single kind of proto substance.Otherwise, the harmonically-proportional expansion of our material world will be impossible, since independent kinds of realities, having the same, similar-equal peculiarities of expansion, seem extremely improbable.By presented judgments, we come to an important conclusion on the single kind of physical reality being the basis of the substance, creating all possible material things.Thus, the next natural question arises: what kind of reality may serve the basis of all?Einstein was deeply convinced; "the electromagnetic field is enough for that" and he stubbornly worked on the idea, until the end of his days (about 30 years), he did not succeed to complete it, however, mostly, because of non-comprehension of the causal essence of gravity as in Ref. [8].
The principal possibility of realization of Einstein's fundamental idea, on a single kind of primordial physical reality, is shown with representation of known elementary particles and their interactions based on the electromagnetic field as in Refs.[1] [4].
It is also possible to mark some observed results in favor of the presented concept.Particularly, the trajectory deflection of NACA's spacecraft "Pioneer" in correspondence to Hubble's constant is possible to observe as non coincidental, as some researchers are inclined to think as Ref. [28].Meanwhile, it can be explained with actual expansion of our Solar system, in accordance to Hubble's universal expansion, with expansion of the material substance in general (that remains invisible to us).
The resent observations of concentrically expanding groups of cosmic objects also may serve as next serious evidence on the expansion concept in general, as in Ref. [29].We see it appropriate referring to a recent publication directly evidencing the rightness of the developed concept of universality of expansion, as Ref. [30].
It is possible to comprehend that universal expansion of the world, consisting of manifold single kind objects will be unobservable (indescribable) within framework of idealized abstract mathematical concepts, without consideration and study of natural properties and peculiarities of physical objects.It means, in complete description of gravity phenomena a couple of known basic natural constants c, h must be expressed.Einstein's GR (and other theories) give us quantitative descriptions of gravity, without causal connection of it to the basic particles of substance that necessary to complete the theory.
Some Quantitative Reasoning on Universal Expansion
We have no intention to put under doubt and challenge the quantitative significance of GR in general, in conformity to previous content.Meanwhile some clarifications, possible simplifications and important conclusions can be easily derived, as we know the causal base of studied phenomena, as we believe.These can be suitable for experimental test.
a) The complicated math apparatus of GR is easy to explain since it relates to "distorted-observable" values and events, and not to "actual events and laws of nature" (1.1.f)), (2.1.5).A second complication with GR is related to universal character of phenomena where "everything" participates in expansion process and all physical values and units become someway-interconnected variables.We can judge from (3.9) that there will be no way to see (measure) the geometrical changes of the expansion process in case of idealized-instant measurements (if we accept c → ∞) in virtue of proportional changes of any real-physical etalons and observation systems in our world.Einstein's equations as well as Schwarzschild solution simply evidence the above-said (3.5), (3.4); all kinds of relative effects become zero if we accept the light velocity as infinite.We will detect action forces between contacting bodies and see "free falling" etc. that will be unexplainable to us, as actually seen.
Some secondary consequences of expansion process only may be detectable by direct vision that becomes possible thanks to limited light velocity.
We see the "far" events, connected to the motion with some delayed time, therefore certain difference relative to events "close" to us can be observed.The light velocity and corresponding parameters of material objects define the significance of "distortion" related to universal expansion that will be arguments in our formulas.Above-said may serve as a causal essence of GR effects in general.
We can illustrate one serious critical remark addressed to GR that puts it under doubt in general, based on above-said, as example.The matter is some theorists have observed that: GR is not adequate from energetic point of view, since the gravitational energy gets different significances depending on the choice of observation systems, as in Ref. [7].
We can comprehend the essence of the problem considering that GR relates to observable values and events, and not to actual ones as said above.Then it becomes clear that-GR is not suitable for causal description of phenomena at all.We must pass into idealized thought system of observation, with absolute constant units of measurement to have the actual picture of gravity phenomena that will correspond to their casual mechanisms, similar to Copernicus' logical operations (1.1.c)).We already used the same operation to explain SR phenomena causally as observable distortions of reality (previous chapter).
The problem, however, becomes technically complicated in this case, because we need to consider new relations, arising due to dependencies of examined physical values on properties and parameters of material substance that we need to consider.b) We will use the local systems of observation, which is maximally comfortable, to simplify our calculus, in comparison to covariant description that used in GR.Moreover, we will examine particular cases only: the homogeneous symmetrical distribution of material substance and the absence of axial-angular momentum also.In this way, we can apply the single coordinate description that changes nothing from conceptual point of view and extremely simplifies the work.c) Initially we intuitively have supposed in the classic physics the opportunity for us to mark absolute static systems and sizes (3.2.f)) and the unchangeable course of time with evenly standard intervals.We have never thought to link our units of measurements with concrete material objects, as we have assume these unchangeable at all, therefore it has not been important how those were set in practice.We can immediately understand now that our "meter" will change proportionally together with all material objects.
However, the question with the "time" is not so easy to solve as we do not have yet an unambiguous definition: what is meant under the term "time" at all! Thus, we are just obligated to ask the natural question: what will happen with our clocks with universal expansion?For this, we must first answer the question, how is "time" linked with material substance?We have used in practice some real physical objects as the etalons of clocks that are able to generate some repeatedly-regular events, the frequency f of which we can adopt as the time course or its inversely-propositional value as the time interval: t = 1/f.We can construct our clock from simplest form of substance-ideal gas.It is possible to conclude that events' frequency in standard condition (for example, the collision of two molecules of gas in certain volume) will be proportional to their distribution density: We adopt (3.10) as definition of "realistic time" connected with substance and variable with the expansion, vs. abstract concept of "time" in supposition of it as unchangeable-absolute.We need to adopt proportionally symmetrical expansion of sizes to all possible material objects that we perceive as "unchangeable" as well: where: V velocity of expansion and R is the distance.The light velocity serves as an important factor that gives us possibility to observe certain secondary effects connected to universal expansion (the GR effects, in analogy to SR effects).Thus, we need to adopt restricted speed of our possible measurements also: where: V ms maximum velocity of measurements (observations).Thus, (3.10), (3.11), (3.12) are the basic principles to describe the expanding world (the consequences of which we perceive as "gravity phenomena").
Let us suppose the expanded substance distributed homogenously by density ρ, in the spherical body with radius R (Figure 5).We can write from (3.10) where: M the mass; k certain coefficient.We can write from (3.11) and from Figure 5.
We define the acceleration of surface point relative to center in radial direction: where: V 0 , R 0 velocity of point on the surface and its radius.
We adopt m i = m g = m (3.8) and a = g 0 because of identity inertia and gravity (3.2.f)) and we get from (3.15) the Newton's law of gravity: where: is the Cavendish constant established experimentally.
We can define the significance V 0 as the final speed in the end of way R from known formula of accelerated movement, accepting R = R 0 : where e V called the escape velocity known within frame of Newton's gravity.We see then V 0 corresponds to v g in (3.4): 0 We can write from (3.17), using the density instead of mass: We can then define k and t from (3.13), (3.15) The significance of t in (3.19) corresponds to virtual time that needs to pass from initial point of expansion (0) to surface, calculated by present scales and units of measurements.The expression (3.19) shows relative meaning of time depending on parameters of material objects.The local character of "time" as property of a concrete material object and wrongness of operation with "universal time" concept become clear by the same.
We emphasize that our definition of "time" (3.10) corresponds to components of stress-energy tensor T μν in Einstein's Equations (3.5) by using known mass-energy relation E = mc 2 : where e ρ is the density of energy.
(3.20) shows that the concept of "time" characterizes the energetic condition of substance.Thus, "time" can have only local-concrete meaning and not abstract universal, separate from material objects that are silently accepted in classic physics.This definition of physical "time" directly corresponds to Einstein's realistic demand (1.1.f)).The relation (3.18) by its form and members corresponds to Hubble's law.We accept the expansion of substance identical to Hubble's expansion (3.The observable closeness ρ → ρ cr may evidence the rightness of this identification.Then the addition of "balancing" constant Λ to Einstein's Equation (3.5) becomes unnecessary, since Hubble's expansion and gravity become the same factor, interconnected and equal (as action with reaction in Newton's 3rd law).We can conclude also that mystically "dark matter" and "dark energy" seems as unnecessary, and the corresponding observable phenomena must be explained in the frame of outlined concept (will discussed).
d) To define relativistic effects connected with universal expansion (gravity) we need to consider 3rd principle (3.12).The interconnected factors of expansion speed and acceleration, in combination with light velocity will be the arguments defining this or that observable relativistic effects, as it derives from above content.We will study the significance of two mentioned factors separately due to comfort of their application for the concrete effects.
Let's assume, an observer measured the distance R 0 using the light signal passing from center to surface (Figure 5).The measured time will be R 0 /c without consideration of expansion.There will be certain delay of time because of expansion speed: We define the relative increase of measured value using (3.22): ( ) Using (3.17), we get: We will be satisfied by examination of "weak" gravity, accepting 0 V c that brings to: The result (3.25) may serve as illustration of physical meaning of Schwarzschild solution and corresponding "curvature of space-time" (3.4).Meantime, we see that (3.24) provides additional correcting members despite incomparable simplicity of used considerations.
We need to consider the increase of expansion speed also during the measurement that will be defined by acceleration factor.The increase of expansion speed will be: We bring causal illustrations to some known GR effects by using deduced factors k v , k g .
"Gravity influence" on the frequency of light
Let us assume light's signal passes way 0 l R near to surface of material body in radial direction.Doppler frequency change will appear in consequence of expansion.Own expansion of l may be ignored in virtue of initial condition and the frequency change will be defined mostly with the acceleration factor.
The summary effect on two opposite radial directions will be defined: The same effect is interpreted with GR as "consequence of difference of gravity potentials": (This effect has been confirmed repeatedly, for example, by Shapiro in laboratory and later by NASA in cosmic scale (Gravity Probe-A).
Deviation of light ray near massive body
The light's ray looks curved in consequence of universal expansion.The massive body M expands from dashed line to solid, accordingly, the observer's position changes during light passage from the edge of material body to the observer (Figure 6).The position of light source seems shifted with angle a.The factors of speed and acceleration of expansion participate in the phenomenon in virtue of scales.We can define the share of expansion speed immediately considering that it is perpendicular to light's ray as bellow: The component of acceleration g n perpendicular to ray only causes the curvature, which is changeable along the way.We introduce variable coordinate x to define g n (Figure 6).
( ) While the ray passes the elementary distance dx the expansion speed will increase in perpendicular direction: We define the summary change of speed accepting that the light's ray passes the distance The resulting deviation angle will be: It coincides with prediction of GR that the first time was confirmed by Eddington in 1919.
The angular displacement of planets' orbits
Causal interpretation of the phenomenon is the same as in above-examined effects.Two factors k v , k g participate in the effect as the observer's location is supposed to be connected to the central body which extremely simplifies the calculations.
The relative expansion of orbit's radius during observation will be: ∆ and angular displacement will be: This prediction of GR also was confirmed by different observations.Meanwhile, one undesirable conclusion follows from a phenomenon that is unspoken.The matter is the light signal "slows down" in case if it passes vs. gravity field (light radiates from central body).The effect gets the opposite sign however, in opposite direction by the same formulas.Then the velocity of light exceeds c that looks contradicting to basic principle of SR!The problem is simple to explain within the expansion concept; the passed way of light's signal l increases (decreases) during the observation (measuring) in consequence of expansion, and the light's velocity does not change.We shall define the time delay considering variability of factors k v , k g in scales of observation 0 l R .
( ) and we get: Similar experiments are also realized by NASA.We shall remark an important point on above examples that all results are based on the same causal concept which may serve as an evidence of its correctness.In author's view: -The possibility of unification of Newton's gravity, Hubble's expansion and Einstein's GR in the same causal context may be adopted as weighty evidence on significance of the offered concept.
The interpretation of other relativistic effects with the same comprehensible causal basis also seems possible.Particularly, the effects connected with presence of angular momentum can be calculated by consideration of the Coriolis Effect arising with expansion process (as Lense-Thirring effect etc.).e) Some "obvious" problems with universal expansion concept arise that need to be examined.
The absence of "gravity waves" follows from offered concept
The question however has following simplest explanation: "gravity waves" have arisen by verbal-arbitrary interpretation of GR, which was illustrated in the previous content.
Einstein's Equations (3.5) are adopted to present the "field equations" and existence of its "indignant states" as well as "gravity radiation" follows formally.The absence of influence on distance, however, is one of the main principles of GR that brought to correct results.Then some internal contradiction arises between "gravity waves" (as the "influence on distance") and "movement by geodetic, free of influence" (the verbal replacement of "field" with the "curvature of space-time" does not remove the problem).
Meantime, it has shown already that GR presents by itself a quantitative description of expanded world and Einstein's equations actually describe the observable movement exclusively.Thus, "gravity wave" becomes a result of misinterpretation of used concepts.The same conclusion derives from direct identification of "gravity" with the "inertia" that is actually adopted in GR (as well as in Newton's gravity).Then it becomes simply obvious that "the inertial waves passing a distance" cannot exist by definition, as the inertial force influence is transferred through direct contact of material objects only.It becomes clear from above-said that experimental detection of "gravity wave" means violation of identity of concepts of gravity and inertia that will be crucial for the offered concept in whole.The experimental confirmation of non-distinguishable gravity and inertia, however, seems easier realizable to us (3.2.f)) than "gravity wave detection" which may remove this problem from the agenda at all, as artificially created.We hope the experimenters may consider the above-said as obvious argumentations and such experiments may be implemented. Problems with "dark matter" and "dark energy" One of the aspects of introduction of "dark matter" was connected to the Hubble's expansion and to the cosmological member, that was examined above (3.3.b), c), d)).Some observed results, however, have pushed to recover the "dark matter" in addition to "dark energy" as well.The issue is, certain cosmic systems have been observed, where significance of Newton's gravity looks as exceeded, while evaluating the quantity of gravitating substance by known criterions.Then theorists introduced some unclear kind of "reality" necessary to com-pensate the "deficit" of gravity.This approach cannot be acceptable to us in virtue of adopted methodology as a pure "ad hoc" hypothesis without any evidence of its existence.Meanwhile, the concept of expansion opens clear opportunity to solve similar cognitive complications as consequences of differences between observable and actual pictures of reality (3.3.a)) (i.e.similar to SR problems).Different "deviations from known natural laws" and causal paradoxes may be observed then that may be pushing to introduce a new hypothetical reality.For example, some increase of planet's orbital movement (3rd point, in 3.3.d)) may be interpreted as "some increase of gravity field" of central body, if desired.The decrease of orbital movement is also possible to observe depending on parameters and observation systems that will demand additional "repulsive forces" with their corresponding sources etc.Then the essence and actual significance of examined problems become clear.We need to notice, that some researchers do not share modern hypotheses on "dark matter"/"energy" and have disputed this approach with certain arguments as in Ref. [31].
Problems with orbital movement and heavenly mechanics
The examined identity of concepts of "gravity" and "inertia" are enough to present the same consequences of gravity as phenomena connected with the universal expansion of substance as our equations remain the same (Newton's gravity as well as GR).However, a huge number of immediate questions beg there with "switch" from gravitation to inertia.The situation is similar to intuitive reaction of the announcement of the Earth as something round, on its movement etc.We will look at couple of probable questions only: -The Moon "feels" through gravity field where the Earth is, and "may choose" the path to move around it.How can it "understand" now: where is the Earth's location?The answer is easy to find considering that expansion is accelerative that is characterized by three physical parameters: direction, speed and acceleration relative to a free, uniform movement that contains in it two parameters only (direction and speed).It follows that the body that is expanding with acceleration, "remembers" the initial point where it has started the process.It will do corresponding oscillations around it by getting certain external impulse.The explanation that the Earth and the Moon have consisted in one common body in the time follows from this.They have been divided by some scenario; both of them now oscillate around the initial common point of their masses, under initial equal impulses on opposite directions.
Next problem concerns the "obvious" conclusion of infinite increase of expansion speed with time due to its accelerating character
A quantitative explanation of the problem is also possible, despite that it sounds somewhat unusual in the framework of adopted traditional concepts.The question is related to "time" concept that was the "universal-abstract" before and now it is directly defined with the density of mass-energy (3.10), (3.20).The inverse proportionality of unit interval of time to density of substance brings to permanent increase of next interval of time in relation to the previous.Thus, permanent increase of speed (the first differential of distance by time) arises as result of decrease of regular events' frequency due to universal expansion.It simply follows from relations (3.13), (3.14), (3.15).The "actual picture" of expansion (that we would see in imaginary absolute system, with unchangeable timers and meters) would be seen by exponential law, where the final speed of expansion strives to the light velocity: V c → < ∞ , with τ → ∞ .i.e., a wrong conclusion arises because we observe a distorted picture of expansion (except Hubble's expansion that we "see" indirectly).We bring some simplest considerations.Using (3.14) and accepting "time" as an independently free variable, we write: . And we get ( ) − .Accepting 0 0 t = we get: ( ) The time change with expansion in imaginary absolute system is illustrated (Figure 7).
Energy of Expansion: Gravity Constant
The gravitational energy is connected to expansion speed, as it follows from the offered concept.It represents kinetic energy, concentrated exclusively in the expanded body.Then its dependence on the system of observation becomes simply explainable.
a) We shall bring here one most important remark for following considerations. From above-said follows, that gravity characteristic of the substance itself has a local significance
It means the peculiarities of expansion are defined by the parameters of a concrete material object.It has been expressed in our formulas, where the adopted "gravity constant" always acts in combination with the massenergy density of substance (or, with "local time" (3.3.c)).Above-said means that "gravity" peculiarity of substance cannot be an independent fundamental constant of nature (in the rank of h, c) and it may be defined within the dependence on certain free parameter of substance, in combination with h, c (3.2.g), 2)).Thus, we must adopt that the experimentally established known to us "gravity constant" becomes variable in the imaginary absolute system, parallel to course of time, with our length units etc. 10 .It becomes clear also, that adopted "gravity constant" is not "so successful" in virtue of possibility of simplest definition, following from (3.14).The relation V/R becomes a sort of local constant for a concrete material object in absolute system of observation, due to the demand of symmetry preservation during expansion (that we observed as Hubble's expansion).Thus, we examine relation V = f(R) for concrete material object to open causal essence of expansion.We shall define kinetic energy of expansion for the standard body M, R 0 relative to its center (Figure 8).
The kinetic energy corresponding to elementary volume between R, R + dR will be: (3.34) appears equal to full "energy of gravitational field" that is calculated by the next imaginary operation: the "gravity source" disintegrates to elementary parts that are shifting away on infinite distance.The resulting energy that needs to be spent to overcome "attraction forces" to complete the described operation gives the same result.Thus, the exact coincidence of results of two considerations confirms the full identity of concepts "gravity" and "inertia" by their energetic significances also (i.e.without "local" restriction).
b) Derivation of the gravity constant: The universal, proportional expansion of material world may be possible in case of a single kind of basic substance, creating all possible kinds of elementary particles, localized (such electron, proton, etc.) as well as non-localized (photon) (3.2.g).3)).The reader can find physical models of elementary particles as localized and non-localized quanta of electromagnetic fields in Refs.[1] [4].Then we can conclude that energy (or velocity) of expansion may have only electromagnetic nature, as it follows from this presentation.We can simply suppose that expansion must be connected to the electromagnetic coupling constant, i.e., it will be defined by the fundamental fine structure constant ( ) as well as all other kinds of interactions of particles and all known physical-chemical peculiarities of substance in general as it presented in Ref. [32].We test the supposition with electron's physical model as localized Compton wave polarized circularly, in view of its simple structure (as wave interference-standing wave) relative to other particles Refs.[1] [4].We will choose certain system of description adopting: [ ] We are using quantum relation (3.9) to define mass of the electron (according to adopted electron model in Ref. [1]).The standing wave and diameter of the particle have some increase in examined model, corresponding to its anomaly magnetic moment: ( )
Conclusions and Discussions
The derived formula for the theoretical significance of gravity constant (3.38) may be interpreted as coinciden-tal, if taken separately.However, in authors' viewpoint, it is difficult to do in the whole context of approach, considering a number of similar "coincidences".Meanwhile, the opportunity of cause-effect, harmonious and self-consistent representation of material world on the unique basis of substance and on the common principles of nature, may be a weighty evidence of significance of the offered concept and used methodology.The productivity of the approach proves the correctness of the wave-dynamical representation of elementary particles and of the microcosm in whole as Ref. [1] [4].It confirms the convictions and demands of undeniable founders Einstein, de Broglie, Schrodinger, Planck and other coryphées of physics unfairly rejected by the majority.
Offered causal interpretation confirms and clarifies quantitative significance of relativity theories; meantime it demands an important revision of the adopted cognitive interpretations (or to replenish their absence).These become conceptually completed, "full-blooded" physical theories, with the same.
New opportunities may open to explain certain problems with cosmology.Particularly, complicated processes of "Gravity collapse" and "Big Bang" may have new aspects of description, remaining in the framework of the known natural laws, i.e. without referring to the mystical "Singularity".
Figure 1 .
Figure 1.P: platform; O: observer; B: bullet; V: speed of bullet; V x : speed of platform; G: the gun; S: the source of light signals; T, M: the clock & meter.
Figure 2 .
Figure 2. S: Light impulse source; T: Timer with photo sencore; l: The rod; M: The mirror.The task of experiment is to define light velocity in moving frame by direct measurements time of passage light pulses of the rod with the length-l.Timer-T starts work when light hits to it from left side of drawing (point 1) and it stops when light hits from right side (3), after reflection (point 2) from mirror: M, fixed in right end of the rod.
Figure 3 .
Figure 3.The difference of trajectories in free falling (a), and in the acceleration (b).
Figure 4 .
Figure 4.The full equality of expansion with the grawity.
Figure 5 .
Figure 5.The expansion of material body.
Figure 7 .
Figure 7. Time change with expansion.
Figure 8 .
Figure 8.The energy of expansion.
where: e λ is the Compton's wavelength for the electron).We write our basic supposition in chosen relative units system in following simplest form: 35) corresponds to escape velocity of electron taking it as the "gravity source".It may be transferred into real units system using below expression: serves as the velocity unit, and the numeric constant e c λ taking into account the difference of local time units in two systems11 (3.3.c)).
λ r is the real-average wavelength as in Ref.[1].We get from (3.37) final expression considering mentioned correction ( ) corresponds to recent measurements. | 26,369 | 2015-09-01T00:00:00.000 | [
"Philosophy"
] |
Amyloid β (Aβ) Peptide Directly Activates Amylin-3 Receptor Subtype by Triggering Multiple Intracellular Signaling Pathways*
Background: Aβ and human amylin peptides share similar biophysical and neurotoxic properties. Results: Aβ directly activates amylin-3 receptor (AMY3) isoform and triggers multiple signaling pathways. Conclusion: Aβ actions are expressed via AMY3 receptors. Significance: AMY3 could serve as a therapeutic target for attenuating Aβ toxicity. The two age-prevalent diseases Alzheimer disease and type 2 diabetes mellitus share many common features including the deposition of amyloidogenic proteins, amyloid β protein (Aβ) and amylin (islet amyloid polypeptide), respectively. Recent evidence suggests that both Aβ and amylin may express their effects through the amylin receptor, although the precise mechanisms for this interaction at a cellular level are unknown. Here, we studied this by generating HEK293 cells with stable expression of an isoform of the amylin receptor family, amylin receptor-3 (AMY3). Aβ1–42 and human amylin (hAmylin) increase cytosolic cAMP and Ca2+, trigger multiple pathways involving the signal transduction mediators protein kinase A, MAPK, Akt, and cFos. Aβ1–42 and hAmylin also induce cell death during exposure for 24–48 h at low micromolar concentrations. In the presence of hAmylin, Aβ1–42 effects on HEK293-AMY3-expressing cells are occluded, suggesting a shared mechanism of action between the two peptides. Amylin receptor antagonist AC253 blocks increases in intracellular Ca2+, activation of protein kinase A, MAPK, Akt, cFos, and cell death, which occur upon AMY3 activation with hAmylin, Aβ1–42, or their co-application. Our data suggest that AMY3 plays an important role by serving as a receptor target for actions Aβ and thus may represent a novel therapeutic target for development of compounds to treat neurodegenerative conditions such as Alzheimer disease.
Alzheimer disease and type 2 diabetes mellitus are age-prevalent diseases that are associated with the deposition of proteinaceous amyloid proteins within the brain (amyloid  protein, A) and pancreatic islet cells (islet amyloid polypeptide, amylin), respectively (1)(2)(3)(4). A and amylin also share biophysical and biochemical characteristics including a propensity to aggregate and form fibrillar structures composed of a core of -pleated sheets (5). Furthermore, amylin inhibits self-association of A into cytotoxic aggregates through direct amylin and A interactions (6,7). A recent proteomics study also indicates that A1-42 and human amylin (hAmylin) 2 share a common pathway for induction of toxicity via mitochondrial dysfunction (8). Also, evidence from pathophysiological, clinical, and epidemiological studies suggests that these two amyloidoses are linked to each other (9 -11). In line with these findings, our studies on cultured human or rat fetal neurons, have shown that electrophysiological and neurotoxic actions of 〈 are strikingly similar to those for hAmylin and that such effects can be blocked by the amylin receptor antagonists AC253 and AC187 (12,13). Additionally, we identified that down-regulation of the AMY3 isoform in neurons using shRNA can blunt the neurotoxic effects of 〈 (12). Collectively, these observations support the notion that A and amylin share a common pathophysiological mechanism possibly via AMY3.
G protein-coupled receptors (GPCRs) constitute a superfamily of cell surface signaling proteins that mediates transduction of a large variety of extracellular stimuli across cell membranes and is involved in numerous key brain neurotransmitter systems, which may be disrupted in pathophysiological conditions such as Alzheimer disease (14). Amylin receptors belong to class B GPCRs and are heterodimeric complexes, formed by calcitonin receptor (CTR) association with receptor activitymodifying proteins (RAMPs) (15). The CTR component of this receptor shares the same general GPCR architecture: seven membrane-spanning interconnected ␣-helices that transmit extracellular signals to the intracellular cytoplasmic signaling cascade. Activation of amylin receptors in mammalian cells has been shown previously to raise cAMP and presumed to involve G S protein (Gs ␣-guanine nucleotide-binding signal transduction protein) and stimulation of adenylate cyclase (16,17).
for the AMY3 isoform of amylin receptors (18). However, the precise intracellular signal transduction pathways following AMY3 activation are not fully understood, and it is not known whether A directly activates AMY3. In the present study, we have for the first time expressed the AMY3 using human embryonic kidney 293 (HEK293) cells to study intracellular signal transduction pathways that are activated by A or hAmylin and to further investigate whether their cytotoxic effects are mediated via this particular isoform of the amylin receptor. (19), respectively. (The pBud-gfp vector was provided by Dr. David Westaway from University of Alberta.) The pBud-gfp-RAMP3 is a bigenic vector, which originally generated from the pBud CE4 vector and contained two distinct promoters upstream for GFP and RAMP3 genes. All of the CTR and RAMP3 gene sequences were confirmed by further DNA sequencing. HEK293 cells were cultured with DMEM (Invitrogen) with 10% FBS (Invitrogen) and grown at 37°C, 5% CO 2 . An AMY3-expressing HEK293 (AMY3-HEK293) cell line was generated by co-transfecting pcDNA3-CTR with pBud-gfp-RAMP3 with a 1:1 molecule ratio using Lipofectamine 2000 transfection reagent (Invitrogen). For controls, AMY1 (CTRϩRAMP1) and AMY2 (CTRϩRAMP2) expressing cells were generated using similar procedures as above. The blank plasmid pcDNA3 and pBud-gfp also were transfected and generated GFP-positive control wild type HEK293 cells. Zeocin TM was used for stable cell line selection and maintenance (200 ng/ml of final concentration). Passages 5-15 of AMY3-HEK293 cells were used for all experiments.
AMY3 Gene Construction and Expression
Immunohistochemistry-For CTR expression, AMY3-HEK293 cells were plated on glass coverslips precoated with poly-L-ornithine (Sigma) in DMEM, 10% FBS, Zeocin medium for 12-24 h. Cells on coverslips were then fixed with 4% paraformaldyhyde and stained with a CTR antibody (12,13,20) (rabbit anti-CTR sera, provided by Dr. P. M. Sexton from Monash University, Victoria, Australia). The secondary antibody was Alexa Fluor 594 donkey anti-rabbit antibody (Invitrogen). Vectashield mounting medium with DAPI (Vector Laboratories Inc., Burlingame, CA) was used to mount slides and DAPI staining. For cAMP determination, mouse monoclonal anti-cAMP (R&D Systems, Minneapolis, MN) antibody was used as a primary antibody and Alexa Fluor 546 goat anti-mouse antibody (Invitrogen) as a secondary antibody. Photomicrographs were imaged using an Axioplan-2 fluorescent microscope with AxioVision software (Carl Zeiss, Toronto, Ontario).
RAMP Gene Expression-To further confirm AMY3 receptor expression, RT-PCR for RAMP1-3 expression was performed.
Initially, RNA was extracted from AMY3-HEK293 and control HEK293 cells using a Qiagen RNeasy kit (Qiagen, Valencia, CA) and following the manufacturer's instructions. RNA samples (500 ng) were reverse-transcribed into cDNA by SuperScript II reverse transcriptase according to the manufacturer's protocol (Invitrogen). One microliter of cDNA was incubated in a 24-l PCR reaction mix (SYBR PCR Mastermix, Qiagen). Primers for human RAMP1, RAMP2, RAMP3, and GAPDH (glyceraldehyde-3-phosphaate dehydrogenase) were purchased from Qiagen (catalog no. PPH02548A, PPH02591A, PPH02536B, and QT01192646, respectively.) Ten microliters of each PCR product were subjected to ethidium bromide gel electrophoresis, and photographs were taken under UV illumination by AlphaImager 2200 (Alpha Innotech, Toronto, Canada).
Quantitative cAMP Measurements-AMY3-HEK293 cells were plated on 24-well plates overnight. Cells were stimulated for 30 min with hAmylin, rat amylin (rAmylin), human calcitonin gene-related peptide, human adrenomedullin, A1-42, and salmon calcitonin 8 -32 over a concentration range (1 pM-1 M). Cellular cAMP levels were measured using a parameter cyclic AMP assay kit (catalog no. KGE002B, R&D Systems) according to the manufacturer's instructions. Data were plotted, and nonlinear regression was fitted with four parameters using Prism software (version 5, GraphPad Software, La Jolla, CA).
Signal Profiling-Intracellular signaling profiles were determined using in-cell Western blot techniques. AMY3-HEK293 cells were seeded at 10,000 cells/well in a 96-well plate (Nalge Nunc Intl., Rochester, NY) in DMEM, 10% FBS, and Zeocin medium. After culturing for 12-16 h, cells were pretreated for 24 h either or not with AC253 (10 M) and then treated with either hAmylin or A1-42 in culture medium for time periods between 10 min and 30 h. Subsequently, cells were fixed with 4% paraformaldehyde for 20 min, permeabilized with 0.2% Triton X-100, blocked with Odyssey blocking buffer (LI-COR, Lincoln, NE), and stained with the following target antibodies. The phospho-p44/42 MAPK (ERK1/2, Thr-202/ Tyr-204), phospho-Akt (Ser-473), and phospho-PKC (II Ser-660) rabbit monoclonal antibodies were purchased from Cell Signaling, Inc. (Danvers, MA). The phospho-PKA R2 (S96) was purchased from Abcam, Inc. (Cambridge, MA), and the cFos rabbit polyclonal antibody was from Santa Cruz Biotechnology (Santa Cruz, CA). The secondary antibody was IRDye 800CW goat anti-rabbit antibody, whereas Sap-phire700 and DRAQ5 were used for cell number normalization (LI-COR). For in-cell Western blot cyclic adenosinemonophosphate (cAMP) quantification, mouse monoclonal anti-cAMP (R&D Systems) was used as a primary antibody, and IRDye 700 goat anti mouse antibody (LI-COR) was used as a secondary antibody. Plates were imaged using an Odyssey Infrared Imaging System (LI-COR), and the integrated intensity was normalized to the total cell number on the same well.
Ca 2ϩ Imaging-Dynamic changes of the free cytosolic Ca 2ϩ concentration were monitored with confocal microscopy as described previously (21). For this, AMY3-HEK293 cells were plated on glass coverslips precoated with poly-L-ornithine and incubated at 37°C for 12-36 h with DMEM, 10% FBS, and Zeocin medium. For Ca 2ϩ imaging, superfusate of ion content similar to extracellular brain fluid thus containing the following: 130 mM NaCl, 4 mM KCl, 1 mM MgCl 2 , 2 mM CaCl 2 , 10 mM HEPES, and 10 mM glucose (pH 7.35) was applied at a rate of 3 ml/min using a roller pump (Watson-Marlow Alitea, Sin-Can, Calgary, AB, Canada). For zero calcium experiment, superfusate as described above except 5 mM, MgCl 2 was used, and no calcium was added. For incubation with the membrane-permeant fluorescent Ca 2ϩ -sensitive dye Fluo-8L-AM (AAT Bioquest, Inc., Sunnyvale, CA), cells were washed twice with superfusate and then incubated with 5 M of the agent for 40 min at room temperature (20 -23°C) within Ͻ2 h before imaging. A1-42 and hAmylin were dissolved in sterile bidistilled water at 1 mM stock solution and incubated at room temperature for 10 min before dilution with superfusate for use at a final concentrations of 0.1-10 M. Fluorescence intensity was monitored with a FV-300 laser-scanning confocal microscope (Olympus FV300, Markham, Ontario, Canada) equipped with an argon laser (488 nm) and excitation/emission filters (FF02-520/28-25; Semrock, Inc.) for an emission wavelength at 514 nm, measured with a numerical aperture of 0.95 20ϫ XLUMPlanF1 objective (Olympus). Images were acquired at scan rates of 1.25-1.43 per second using a 2-3ϫ digital zoom at full frame (512 ϫ 512 pixel) resolution. Regions of interest were drawn around distinct cell bodies, and analysis of time courses of change in fluorescence intensity were generated with FluoView software (version 4.3; Olympus).
MTT Cell Death Assay-AMY3-HEK293 cells were seeded to 5000 cells/well in a 96-well plate in DMEM, 10% FBS, and Zeocin medium and incubated overnight. Cells in culture medium were preincubated for 24 h either with or without AC253, KH7 (selective soluble adenylyl cyclase inhibitor, R&D Systems), or FR180204 (selective ERK inhibitor, R&D Systems) and followed by treatment with either hAmylin, A1-42, or A42-1 for 24 -48 h. At the end of treatment, 20 l of 5 mg/ml methylthiazolyldiphenyl-tetrazolium bromide (MTT; Sigma) was added to each well and incubated at 37°C for 3 h. Medium was removed, 100 l of MTT solvent (isopropanol with 4 mM HCl) added to each well, and the plates were incubated for 30 min at room temperature on a rotating shaker. Plates were analyzed on a microplate reader at a 562-nm wavelength.
Drug Treatment-Cell dysfunction and cell death resulting from exposure to A or hAmylin are mediated by soluble small or intermediate oligomers, whereas large, insoluble deposits might function as reservoirs of the bioactive oligomers (5,22,23). Soluble oligomeric A1-42, the reverse non-functional sequence peptide A42-1, hAmylin were used in the present study and were prepared here as per published protocols (12,24). Human variants of CGRP and adrenomedullin), as well as rAmylin, and salmon calcitonin (8 -32), which are also members of the same CGRP family, were used here for a comparison with the effects of A1-42, A42-1, hAmylin, and AC253. AC253 (LGRLSQELHRLQTYPRTNTGSNTY) is a polypeptide that is a potent amylin receptor antagonist similar to AC187 or salmon calcitonin (8 -32), which also display antagonist activity at the amylin receptor (12,13,25). A1-42 and A42-1 were purchased from rPeptide (Bog-art, GA), and hAmylin and AC253 were purchased from American Peptide (Sunnyvale, CA).
Statistical Analysis-Values are expressed as mean Ϯ S.E. Statistical analysis was performed by one-way analysis of vari- ance followed by Tukey's test when appropriate using Prism software. p Ͻ 0.05 was taken as significant.
RESULTS
Stable Expression of AMY3 in HEK293 Cells-We first generated the novel, AMY3-HEK293 stable cell line using co-trans-fection of the genes for the AMY3 constituents, CTR, and RAMP3 (Fig. 1a), in a 1:1 molecule ratio with the assumption of equal CTR and RAMP3 expression in individual cells. As an indication of RAMP3 expression, 60 -70% of AMY3-HEK293 cells were GFP-positive 24 h following co-transfection, whereas Ͼ95% of cells were GFP-positive already after 3 passages with , n ϭ 4). e, shown are changes in cAMP after exposure of AMY3-HEK293 cells to a full range and concentrations of different peptides derived from the calcitonin gene-related peptide family. The range of EC 50 values for A1-42, hAmylin, and rAmylin are 7.7 (6.3-9.6), 2.4 (1.7-3.6), and 6.9 (5.0 -3.6) nM, respectively. Data were plotted and nonlinear regression fitted with four parameters using Prism software; data are from six wells at each concentration of the three peptides. sCal 8 -32 , salmon calcitonin 8 -32. zeocin selection using 200 ng/ml. RT-PCR further confirmed RAMP3 gene expression in the AMY3-HEK293 cells. There is weak endogenous RAMP1 and RAMP2 gene expression in HEK293 cells that does not change after AMY3 expression. There is also little endogenous RAMP3 gene expression in HEK293 cells (Fig. 1b). Fig. 1c shows GFP-positive cells are all stained with CTR, which indicated successful AMY3 expression.
AMY3 Activation with A1-42 and hAmylin Increases Cellular cAMP-As indication that the AMY3 is functional in the new cell line, we first identified that cAMP increases observed after 30 min of exposure to the established agonist for this receptor, hAmylin (1 M), could also be evoked following application of A1-42 (1 M) (Fig. 2a). As expected, hAmylin and A1-42 did not evoke cAMP increases in control wild type cells (Fig. 2b). There also was no significant increase in cAMP following exposure to hAmylin or A1-42 in HEK293 cells that expressed single components of AMY3, i.e. either a functional CTR (supplemental Fig. 1) or RAMP3. Next, we quantified the effects of hAmylin (0.1-10 M) on cAMP levels with in-cell Western blot. In fact, both agents increased cAMP in a very similar, concentration-dependent manner (Fig. 2c). In contrast, cAMP was not affected by A42-1, human calcitonin gene- related peptide or human adrenomedullin (Fig. 2d). To further validate functional expression of AMY3 in the HEK cells, a cAMP assay over a full concentration range of the different peptides was performed. A1-42, hAmylin, and rAmylin concentration-response curves were non-linearly fitted (Fig. 2e). The EC 50 values for A1-42, hAmylin, and rAmylin are 7.7, 2.4, and 6.9 nM, respectively.
AMY3 Activation with A1-42 and hAmylin Increases Cytosolic Ca 2ϩ -Confocal microscopy was used to investigate whether signaling pathways in AMY3 activation include the important intracellular second messenger Ca 2ϩ (26). Under control conditions, fluorescence intensity in cells loaded with the fluorescent Ca 2ϩ dye Fluo-8L-AM did not show notable spontaneous fluctuations of cytosolic Ca 2ϩ (Fig. 3a). Bath application of hAmylin (0.1-2 M) for 30 s produced a major Ca 2ϩ increase within Ͻ1 min after entry of the peptide within the imaging chamber. These Ca 2ϩ increases displayed a sharp peak, indicating that return to base line already started at the end of the application period of the peptide or within Ͻ30 s after start of return to control superfusate. Recovery to base line was achieved typically within Ͻ2 min after return to control perfusing solution (Fig. 3b). Also A1-42 (0.1-2 M) increased cytosolic Ca 2ϩ in a fashion very similar to that observed for hAmylin (Fig. 3c). In contrast, Ca 2ϩ did not rise in response to hAmylin (0.5 M, 30 s) following superfusion of cells with AC253 (2 M) for 30 s prior to application of the agent, whereas AC253 alone (1-10 M) also did not change Ca 2ϩ base line (Fig. 3, e and f). Similarly, the Ca 2ϩ increase due to A1-42 (0.5 M, 30 s) was abolished by preincubation of AC253 (2 M) (Fig. 3g). Co-application of A1-42 (0.25 M) with hAmylin (0.25 M) elevated Ca 2ϩ levels to a similar extent as when each drug was applied alone, and this response to a combined application of the two peptides also was abolished by 2 M AC253 (Fig. 3, d and h). Bar graphs in Fig. 3i show quantification of these data. Fig. 3j and the accompanying videos (supplemental data) show time-lapsed recordings of Ca 2ϩ signals from the same AMY3-HEK293 cell in response to hAmylin, A1-42, hAmylinϩ A1-42, AC253, AC253ϩhAmylinϩA1-42, and recovery. In wild type (non-transfected) HEK293 cells, A1-42 or hAmylin did not induce increases in cytosolic Ca 2ϩ (data not shown). Increases in cytosolic Ca 2ϩ levels after AMY3 activation with hAmylin and A1-42 mainly depend on extracellular calcium. In the presence of calcium-free superperfusate, hAmylin and A1-42 only produced small and delayed cytosolic Ca 2ϩ increases (ϳ20% of peak increases of cytosolic Ca 2ϩ under normal calcium concentration, Fig. 4).
Signaling Pathways Involved in A1-42 Activation of AMY3-In addition to their effects on cAMP and Ca 2ϩ , we found here that A1-42 and hAmylin also increase phosphorylation of the type II subunit of cAMP-dependent protein kinase A (PKA R2) (Fig. 5a), activate mitogen-activated protein kinase (ERK1/2) (Fig. 5b), increase cellular levels of the transcription factor cFos (Fig. 5c), and increase phosphorylation of Akt (protein kinase B (pAkt)) (Fig. 5d). All of these effects were blocked by AC253 (Fig. 5, a-d) and not mimicked by A42-1 (data not shown). In another series of experiments, phosphorylation of protein kinase C was not observed to change with exposure to either A1-42 or hAmylin (data not shown).
In a further set of experiments, we elucidated that phosphorylation of ERK1/2 or PKA depends on the duration of exposure to hAmylin or A1-42. The phosphoERK1/2 peak is reached ϳ10 min after start of application and returns to base- line level within 0.5 h (or slightly more or less) after exposure to either hAmylin (0.2 M) or A1-42 (0.5 M) (Fig. 6a). The phospho-PKA R2 also reaches a maximum level at ϳ10 min, but the phosphorylated form of this protein lasts longer and returns to base-line level at 2 h (or slightly more or less) for A and hAmylin (Fig. 6b).
A1-42 Activation of AMY3 Triggers Cell Death at Higher Concentrations and Longer
Exposure Times-Our previous observations that 〈1-42 and hAmylin induce apoptotic cell death in cultured neurons suggested that these effects may require AMY3 activation (12,13). This hypothesis is supported by our present findings using an MTT assay that 48 h incubation with either A1-42 (2-20 M) or hAmylin (0.2-2 M) causes cell death of AMY3-HEK293 cells in a concentration-dependent manner (Fig. 7a). Next, we confirmed that hAmylin and A1-42 induced cell death is dependent on AMY3 expression. In cells expressing other amylin receptor subtypes (AMY1 or AMY2, supplemental Fig. 2) or single components of AMY3 (CTR or RAMP3), hAmylin, and A1-42 did not induce significant cell death (Fig. 7b). The cell death observed when AMY3-HEK cells were exposed to hAmylin or A1-42 occurred after 24 h (Fig. 7c). Furthermore, cell death induced from hAmylin and A1-42 is attenuated significantly by pretreatment with AC253 (Fig. 7d). Interference with downstream mediators of AMY3 activation using KH7 (adenylate cyclase inhibitor) or FR180204 (an ERK1/2 inhibitor) also protected AMY3 cells from hAmylin or A1-42 cytotoxicity (Fig. 7d).
DISCUSSION
Our data demonstrate that both A1-42 and hAmylin act as agonists at the AMY3 subtype that has been expressed here for the first time using HEK293 cells. Amylin receptors are heterodimerized by CTR and one of three RAMPs, thus generating multiple amylin receptor subtypes, AMY1-3 (Fig. 1a) (15). There are several CTR isoforms in the human (27). In this study, hCTRa was used for AMY3 receptor construction, which is insert-negative and modulates cell cycle progression (28). The importance of receptor splice variation in AMY physiology remains to be elucidated. These amylin receptor subtypes are pharmacologically distinct and demonstrate different binding affinities to members of the calcitonin peptide family, which includes CGRP, adrenomedullin, and amylin (15). However, due to the lack of selective pharmacological tools, significant complexity of this system and lack of specific RAMP antibodies, it has been difficult to confidently assign specific amylin functions to one of these receptor subtypes. Data from functional bioassays, including binding studies (16,29,30), demonstrate that AMY3, which is a heterodimeric complex of CTR and RAMP3, has a high affinity for amylin. Amylin receptor components (CTR and RAMPs) are reported to be distributed widely in the central nervous system with pronounced expression within spinal cord, brain stem, cortex, hypothalamus, and hippocampus (31,32). However, data on distribution of specific subtypes of RAMPs within the CNS is lacking. Moreover, at present, there is no information on the functional effects of co-localization of CTR with RAMP3 to generate AMY3. Herein, we provide evidence that the AMY3 subtype is indeed the specific target receptor for direct actions of A (and hAmylin) at the level of the cellular membrane.
For the cAMP production assay, the mean EC 50 for hAmylin (2.4 nM) in our AMY3 cell line is close to a previous reported value for rAmylin (18). We have identified that AMY3 activation results in G␣ s -mediated adenylate cyclase activation, with a subsequent increase in cAMP and activation of PKA. This occurred in a manner similar to that reported for the CGRP receptor, which is also a member of the same family of CTR (33). PKA is a multiunit protein kinase that mediates signal transduction of GPCRs following their activation by adenyl cyclase-mediated cAMP formation and is involved in a wide range of cellular processes. PKA R2 is one of the regulatory isoforms of the enzyme, which is predominantly expressed in adipose tissue and brain (34). Nearly all PKA activity in adipose tissue and 50% of PKA activity in the striatum, hypothalamus, and cortex is attributed to the subunit. Disruption of PKA R2 affects physiological mechanisms known to be associated with healthy aging in mammals, which include increased lifespan and decreased incidence and severity of a number of age-related diseases in PKA R2 null mice. In that context, our data further indicate that PKA R2 is secondarily activated after A or hAmylin stimulation of AMY3 and that such an effect may contribute to the long term deleterious neuronal actions of A or amylin as indicated by the present findings of increased cell death of AMY3-HEK293 cells exposed to these peptides.
The increases in cytosolic Ca 2ϩ observed here following hAmylin or A1-42 application could also occur via G␣ q activation. G␣ q proteins activate phospholipase C, which cleaves phosphatidylinositol-4,5-bisphosphate into diacyl-glycerol and inositol trisphosphate, leading to mobilization of Ca 2ϩ from cellular stores (35). Ca 2ϩ represents a ubiquitous intracellular second messenger with enormous versatility (26). The versatility of Ca 2ϩ as a signaling molecule is based on its binding kinetics, varying amplitude, spatiotemporal distribution, and ability to cross-talk with multiple other signaling cascades within the GPCR-activated pathways, including ERK and Akt. The Akt serine/threonine kinase (also called protein kinase B) has emerged as a critical signaling molecule within eukaryotic cells and regulates diverse aspects of neuronal cell function protein translation and cell size, axonal outgrowth, suppression of apoptosis, and synaptic plasticity. Akt activation could result from activation of phosphatidylinositol 3-kinase (PI3K) or G␣ q activation.
Inhibition of adenylyl cyclase at a lower concentrations of KH7 (1-2 M) can protect AMY3 cells from hAmylin and A1-42 damage. However, at higher concentrations (4 M), KH7 did not demonstrate a protective effect, which could be related to interruption of normal cellular function at such concentrations. Altering activity of adenyl cyclase results in changes in cAMP second messenger levels, which in turn affects PKA and protein phosphatase A activity. Protein phosphatase A is inhibited by increased PKA activity, thus maintaining Akt in an activated state (36). Both hAmylin and A1-42 activate Akt, which indicates that AMY3 may play also an important role in controlling cell fate. Most likely, changes in Akt activity that we observed are secondary to alterations in cAMP and PKA that result from AMY3 activation. The rapid Ca 2ϩ increase associated with AMY3 activation could also contribute to Akt activation. ERK1/2 is a member of the mitogen-activated protein kinase family that is centrally involved in many processes during the lifetime of a cell. This kinase has been not only associated with proliferation, differentiation, and protection against apoptosis but also has been linked to cell death (37). A selective ERK1/2 inhibitor, FR180204, can protect AMY3 against hAmylin or A1-42 toxicity, but this protective effect does not appear to be concentration dependent. The magnitude and the duration of ERK1/2 activity may determine its role in regulating different aspects of cellular function. Following AMY3 activation, temporal changes in ERK1/2 phosphorylation followed alterations in cAMP and PKA activity. At early stages or with short term activation, ERK1/2 is associated with cell proliferation, regulation of cellular function, and increases in the transcription factor, cFos. However, after longer periods of activation, ERK1/2 appears to trigger cell death pathways, an observation that seems consistent with our finding of phosphorylation of this kinase and the resultant cell death as shown by our MTT assay data.
In conclusion, we provide for the first time evidence that A1-42 directly activates AMY3 and triggers several intracellular signaling pathways, including cytosolic cAMP and Ca 2ϩ rises. AMY3 likely regulates cellular functions by changing activity of PKA, ERK1/2, and Akt. The sustained activation of AMY3 triggers phosphorylation of ERK1/2 resulting in cell death. Putative uncontrolled elevations of cytosolic Ca 2ϩ as a result of prolonged AMY3 activation may also perturb homeostasis of the endoplasmic reticulum, produce mitochondrial dysregulation and engagement of caspases that contribute to apoptosis. The possible pathophysiological mechanisms whereby A1-42 activates AMY3 and triggers multiple signaling pathways are illustrated in Fig. 8. Our data suggest the AMY3 is receptor target for the actions of human amylin and A and may play an important role in pathogenesis of conditions, where these amyloidogenic proteins have been implicated, namely type 2 diabetes mellitus and Alzheimer disease. Thus, it may be possible, for example, to develop novel therapies for Alzheimer disease by altering AMY3 function or its downstream signaling pathways through the design of highly selective antagonists for this isoform of amylin receptor family. | 6,158.6 | 2012-04-12T00:00:00.000 | [
"Biology",
"Chemistry"
] |
Self-Tracking Emotional States through Social Media Mobile Photography
This paper presents a preliminary breakdown of the results obtained in an exploratory study conducted through the mobile application Instagram. Our goal was to inspect the potential benefits of combining the self-reporting of emotions with everyday mobile photographic practices to learn more about users’ experiences. To do so, we instructed 25 participants to assess and report their emotional states using Instagram, during a 4-week period, according to a pre-established set of instructions. Participants also filled in pre and post-study questionnaires. We then analysed the 291 submissions obtained and the results from the questionnaires, focusing on three aspects: the categories of the photographs taken, the emotion labels used, and the feedback provided by the participants. We end by presenting and discussing a set of insights that might be useful in the design of mobile apps to improve emotional self-awareness and wellbeing.
INTRODUCTION
The self-tracking of emotional states (also known as Mood Tracking) consists in habitually monitoring our emotional states, over extended periods of time, to better understand how, and why, they vary.With that information in hand, we can better learn how to improve our emotional wellbeing.The selftracking of emotions is a typical exercise within the clinical practice in mental health, where therapists often ask their patients to keep mood diaries.This exercise can be beneficial, not only for those experiencing depression, bipolar disorder, and other dysfunctions; but also for those who wish to know more about themselves.In fact, mood tracking is becoming increasingly prevalent among the general public thanks to the myriads of applications that exist for that purpose.The process of self-tracking has two central aspects: collection (of relevant personal information) and reflection (to produce insights) [13].The work presented here focuses on the collection (selfreporting) of one's emotional states.
LITERATURE REVIEW
Self-reporting one's affective states is the most traditional method used to collect a person's emotional data within research [12].Still, a systematic literature review [7] established that the self-report of emotional states was under-researched and that additional studies involving users in real contexts and interfaces with new interaction styles are needed.
Unfortunately, self-tracking technologies often present issues of low adherence and high dropout rates due to barriers like lack of motivation, lack of time, or no immediate access to the tracking tool [9,13,23].Some social media mobile apps, however, seem to be able to overcome these barriers.Nowadays, people usually carry their mobile phones with them all the time and recurrently use apps like Instagram, which counts with more than 700 millions users as of April 2017 1 .For many of those users, sharing photos of their lives, along with their thoughts and feelings, has become a vital part of their daily routines, making social media a naturalistic setting, laden with emotional content.Another social media app, Facebook, lets its users add their current emotional state when they make a post (e.g., I am feeling angry) and to express their feelings about others' posts through "reactions" like "Love" or "Sad".Plus, on Twitter and Instagram individuals often use hashtags to express their emotions (e.g., #bored).There are many studies about emotions and social media, and while most focus on sentiment analysis -"the automatic extraction of sentiment-related information from text" [24], there is also research focused on directly exploring the potential of social media to help better understand mental health and depression [1,2,5,14,16], wellbeing [4] and other relevant public health matters [8].
We decided to use the application Instagram in our research because it offers a type of data that is particularly interesting in the study of behaviours related to emotions: photographs.The richness of photographs allows individuals to communicate messages in a way that would not be possible only with words [1].
Participants
To obtain the corpus of data that supports this paper, we ran a study with 25 participants (17 female), with ages between 25 and 63 years old (mean=34.4),over the course of 4 weeks.We gathered 291 submissions, with an average of 11.2 submissions per user.We recruited the participants via announcements placed on social networks, and we ran a prize draw for a 50€ voucher at the end of the study, as an incentive.All participants were informed about the procedures and agreed to a consent statement.We also asked for permission to use their photographs in future scientific publications.Additionally, we asked participants if they preferred to share their photos openly on Instagram or submit them privately, through the app's direct messaging system -3 participants chose the latter.Most participants (70.8%) used Instagram regularly before the study.
Procedure
We began the study with a questionnaire, meant to obtain demographic data and gauge the participants' habits regarding their wellbeing.A 4week data-collection period ensued, followed by another questionnaire, designed to get feedback on the participants' overall experience throughout the study.During the data-collection period, we asked the participants to share at least two photographs per week on their Instagram account, according to specific instructions, which were detailed in a document and summarized on a cheat sheet.Figure 1 exemplifies the application of these instructions (succinctly described in the list below).The elements of the instructions that directly concern the self-reporting of emotions are the emotion labels and two well-established dimensions of emotional responses: arousal and valence [20,21], which we decided to rename as energy and pleasantness, respectively, as we considered these to be simpler for the participants to understand and remember.Valence (pleasantness) describes the degree to which an emotion is either positive or negative, and arousal (energy) is the intensity of the felt emotion.We believe that understanding emotional states through these two dimensions can be very enlightening.Assessing the degree of arousal points to the vital relationship between emotions and body activation (more specifically, autonomic nervous system activity).Understating how emotions manifest themselves through physical cues in our bodies is vital to increase emotional self-awareness [6].Labelling emotional states and improving our emotional vocabulary is also essential to better recognize and communicate what we are feeling.The lists of emotions and colours initially used in the instructions stem from Plutchik's Psychoevolutionary Theory [17].We chose this particular theory because the list of (primary) emotions was short (8) and had an established relationship between colours and emotions.However, during the second week, some participants started to complain about the complexity of the tagging process.Thus, we decided to revise the instructions, at the beginning of the third week, with 72% of the participants adhering to the new rules.The adjustment consisted of letting the users freely add whatever emotions and colours they wanted, using hashtags (e.g., #happiness, #lilac).We were aware that this modification would complicate the data analysis, but since the study had an exploratory nature we decided to proceed with the adjustment.
DATA ANALYSIS
Since we received some of the entries via direct message, and some participants did not have a public profile on Instagram, we could not pull the data from Instagram's API.Instead, we collected it manually, by iteratively going through each participant's profile and downloading the photographs and transcribing the corresponding text.We then categorized the 291 photographs obtained and assessed the inter-rater reliability (IRR).We also examined all the emotion labels used in the entries and the feedback given by the participants in the final questionnaire.
Photographs
One of the authors (coder 1) used a thematic analysis procedure [3] to systematically find the underlying themes within the data, which resulted in a set of seven categories.Then, the remaining authors (coders 2 and 3) independently coded the pool of photos into the seven pre-established categories.The two authors did not receive any instructions other than the categories names.We computed the Cohen's Kappa for each pair of coders.The first pair (coder 1 and coder 2) obtained a very good strength of agreement (κ = 0.819).The second pair (coder 1 and coder 3) obtained a good strength of agreement (κ = 0.790).These results demonstrate consistency among the observational ratings given by the coders.The seven visual themes are Nature, Surroundings, Objects, People, Food, Animals and Pictorial Images.Figures 2 to 8 show some photos from the study from each category.There are some parallelisms between these categories and the ones reported in a recent study where 1000 photographs obtained through Instagram's API were methodically analysed [10].In that study, there was also a category labelled Food, two categories related to People (Friends and Selfie), one category similar to Animals (Pets), and two categories that could be potentially included in our Objects category (Gadgets and Fashion).
Emotion Labels
The chart in Figure 9 shows the number of posts that were tagged with each emotion label.The chart includes all the labels that appeared in at least four submissions.The label "joy" was by far the most employed label, appearing in a total of 159 posts (56,58%).Plus, if we group the label "joy" with the similar label "happy", we have a total of 175 posts (62,27%).The second most-used label was "anticipation", present in 64 submissions (22,77%), followed by "trust" and "surprise", with 26 posts each (9,25%).The label "sadness" was present in only 19 posts (6,76%).This data is consistent with what is described in the literature: individuals are more disposed to share positive events than negative ones, due to stigma and selfpresentation concerns [1].Furthermore, sharing only favourable and socially desirable images is associated with positive outcomes like higher self-esteem, since it helps people affirm positive views of themselves [1,25].
Number of posts
One participant said "I admit I was auto-censored to upload pictures in "bad days" because I'm not a person who wants to expose my problems to the world.",moreover, another stated, "I realized that in fact the tendency is to put into Instagram things that make me feel positive emotions.".
Participants' Feedback
According to the feedback provided in the final questionnaire, participants enjoyed the overall experience of assessing and reporting their emotional states while taking photographs with their mobile phones.From their comments (Table 1), it transpires that they found the study beneficial, especially regarding the self-assessment of emotional states.They reported that it made them reflect more and be further aware of their emotions."This is a good exercise to make us be more aware of our emotions, and I think that's half way for a healthy mind." "It was useful as a self-assessment and reflection about (what) I was doing." To the question about whether or not the emotion assessment made using the pleasantness (valence) and energy (arousal) dimensions was easy to perform, 88% of the participants answered positively.The remaining 12% said that the two dimensions were difficult to quantify.It was interesting to notice that some participants seemed to believe that the two dimensions always had to go in tandem (e.g., one could only feel high energy when pleasantness was high as well, and vice versa).One of the participants stated "I guess that being aware of my energy and pleasantness levels made me do something about that, especially when they were low.", while another claimed "I was surprised how sometimes, we can be overly tired but completely satisfied with the situation.",referring to a high pleasantness and low energy situation.This mental model is incorrect since the two dimensions can have opposite values.For instance, distress is characterized by high energy and low pleasantness, whereas contentment is its bipolar opposite [20].One participant commented on how "it became clear that energy and pleasantness aren't always aligned."
DISCUSSION
After reviewing the data collected during the study, we are inclined to believe that combining social mobile photography with the self-reporting of emotions is a viable way of incorporating this practice into people's everyday lives, as a means to fight issues of attrition commonly present in this type of technologies.We learned that people are generally not disposed to share their negative emotions with others, and thus it is important to offer users the possibility to keep entries private in a future mobile app so that they can feel at ease to explore their negative emotions as well.We also believe that it is imperative to better explain and illustrate the concepts of arousal and valence before asking users to deal with these dimensions.
Explaining the concepts better in the instructions and reporting interface, and experimenting with other labels, such as "body activation" instead of energy or arousal, might help.Lastly, we trust it is valuable to be aware of the categories of photographs that people usually take and share when designing and developing apps, and even perhaps use photos from the categories during the design process (e.g., mockups, personas, user stories).
CONCLUSIONS AND FUTURE WORK
This paper presented the first impressions of an exploratory, 4-week long study, where 25 participants submitted, through the mobile application Instagram, a total of 291 photographs and accompanying textual data regarding the selfreporting of emotional states.The primary goal of this experiment was to learn more about the process of self-assessing and reporting emotional states while taking photographs with a mobile app, and we intend to further explore this matter with further analysis of the data, interviews and participatory design sessions.We also intend to perform thorough individual analyses of each participant's data to understand the users better and produce reports to investigate how we can improve the design for the reflection aspect of the self-reporting of emotional states.Furthermore, we expect to use some of the gathered data, like photographs and captions, in the sketching process of a future mobile app for emotional wellbeing.Finally, thanks to the positive comments and reactions of the participants, we feel encouraged to keep exploring the frontiers between emotions, mobile phones, and photographs, to inform the design of future mobile technologies to promote emotional awareness and wellbeing.
Figure 1 :
Figure 1: A submission made by on the participants.
Figure 2 :
Figure 2: Set of photographs from the Nature category.
Figure 3 :
Figure 3: Set of photographs from the Surroundings category.
Figure 4 :
Figure 4: Set of photographs from the Objects category.
Figure 5 :
Figure 5: Set of photographs from the People category.
Figure 6 :
Figure 6: Set of photographs from the Food category.
Figure 7 :
Figure 7: Set of photographs from the Animals category.
Figure 8 :
Figure 8: Set of photographs from the Pictorial Images category.
Figure 9 :
Figure 9: Chart showing the number of posts tagged with each emotion label. | 3,278.8 | 2018-07-01T00:00:00.000 | [
"Computer Science",
"Psychology"
] |
Quadruple Integral Involving the Logarithm and Product of Bessel Functions Expressed in Terms of the Lerch Function
In this paper, we have derived and evaluated a quadruple integral whose kernel involves the logarithm and product of Bessel functions of the first kind. A new quadruple integral representation of Catalan’s G and Apéry’s ζ(3) constants are produced. Some special cases of the result in terms of fundamental constants are evaluated. All the results in this work are new.
Significance Statement
Bessel functions were first studied by Daniel Bernoulli [1] and then generalized by Friedrich Bessel [2] and are canonical solutions of Bessel's differential equation (see section (10.13) in [3]). Bessel functions are often used as approximants in the construction of uniform asymptotic approximations and expansions for solutions of linear second-order differential equations containing a parameter (see section (10.72) in [3]). Bessel functions are also used in the physical problem involving small oscillations of a uniform heavy flexible chain (see section (10.73) in [3]). Bessel functions arise in the application of cylindrical symmetry in which the physics is described by Laplace's equation (see section (10.73) in [3]).
The definite integral of the product of Bessel functions, which find importance in many branches of mathematical physics, elasticity, potential theory and applied probability, is studied in the works of Glasser [4] and Chaudhry et al. [5]. Multiple integrals of Bessel functions are used in the geometry of fractal sets and studied in the works of Falconer [6] and Ragab [7].
In this work, our goal is to expand upon the current literature of multiple integrals involving the product of Bessel functions by providing a formal derivation in terms of the Lerch function. It is our hope that researchers will find this new integral formula useful for current and future research work where applicable. Consequently, any new result on multiple integrals of the product of Bessel functions is important because of their many applications in applied and pure mathematics.
Introduction
In this paper, we derive the quadruple definite integral given by where the parameters k, a, b, p, q, v and m are general complex numbers. The derivations follow the method used by us in [8]. This method involves using a form of the generalized Cauchy's integral formula given by where C is in general an open contour in the complex plane, where the bilinear concomitant has the same value at the end points of the contour. We then multiply both sides by a function of x, y, z and r, and then take a definite quadruple integral of both sides. This yields a definite integral in terms of a contour integral. Then, we multiply both sides of Equation (2) by another function of x, y, z and r and take the infinite sum of both sides such that the contour integrals of both equations are the same.
Definite Integral of the Contour Integral
We use the method in [8]. The variable of integration in the contour integral is α = w + m. The cut and contour are in the first quadrant of the complex α-plane. The cut approaches the origin from the interior of the first quadrant, and the contour goes round the origin with zero radius and is on opposite sides of the cut. Using a generalization of Cauchy's integral formula, we form the quadruple integral by replacing y by log ayz rx and multiplying by from Equation (3.326.2) in [9] and Equation (521.1) in [10], where −1/2 < Re(m) < Re(v) + 1, −1/2 < Re(w + m), Re(b, c) > 0, and using the reflection Formula (8.334.3) in [9] for the Gamma function. We are able to switch the order of integration over α, x, y, z and r using Fubini's theorem since the integrand is of bounded measure over the space
The Lerch Function and Infinite Sum of the Contour Integral
In this section, we use Equation (2) to derive the contour integral representations for the Lerch function.
(r m x m y n z n − y m z m r n x n )dxdydzdr Proof. Use Equation (7) and form a second equation by replacing m → n and take their difference. Next, set k = −1, a = 1, b = c = 1 and simplify using entry (2) in the table below (64:12:7) in [11].
Discussion
In this paper, we have presented a novel method for deriving a new quadruple integral involving the product of Bessel functions along with some interesting special cases using contour integration. We will use our method to expand upon this current work and derive other multiple integrals involving other special functions. The results presented were numerically verified for both real and imaginary and complex values of the parameters in the integrals using Mathematica by Wolfram.
Author Contributions: Conceptualization, R.R.; methodology, R.R.; writing-original draft preparation, R.R.; writing-review and editing, R.R. and A.S.; funding acquisition, A.S. All authors have read and agreed to the published version of the manuscript. | 1,157.8 | 2021-11-30T00:00:00.000 | [
"Mathematics"
] |
Detection prospects for the second-order weak decays of $^{124}$Xe in multi-tonne xenon time projection chambers
We investigate the detection prospects for two-neutrino and neutrinoless second order weak decays of $^{124}$Xe -- double electron capture ($0/2\nu\text{ECEC}$), electron capture with positron emission ($0/2\nu\text{EC}\beta^+$) and double-positron emission ($0/2\nu\beta^+\beta^+$) -- in multi-tonne xenon time projection chambers. We simulate the decays in a liquid xenon medium and develop a reconstruction algorithm which uses the multi-particle coincidence in these decays to separate signal from background. This is used to compute the expected detection efficiencies as a function of position resolution and energy threshold for planned experiments. In addition, we consider an exhaustive list of possible background sources and find that they are either negligible in rate or can be greatly reduced using our topological reconstruction criteria. In particular, we draw two conclusions: First, with a half-life of $T_{1/2}^{2\nu\text{EC}\beta^+} = (1.7 \pm 0.6)\cdot 10^{23}\,\text{yr}$, the $2\nu\text{EC}\beta^+$ decay of $^{124}$Xe will likely be detected in upcoming Dark Matter experiments (e.g. XENONnT or LZ), and their major background will be from gamma rays in the detector construction materials. Second, searches for the $0\nu\text{EC}\beta^+$ decay mode are likely to be background-free, and new parameter space may be within the reach. To this end we investigate three different scenarios of existing experimental constraints on the effective neutrino mass. The necessary 500 kg-year exposure of $^{124}$Xe could be achieved by the baseline design of the DARWIN observatory, or by extracting and using the $^{124}$Xe from the tailings of the nEXO experiment. We demonstrate how a combination of $^{124}$Xe results with those from $0\nu\beta^-\beta^-$ searches in $^{136}$Xe could help to identify the neutrinoless decay mechanism.
I. INTRODUCTION
The origin of the matter-antimatter asymmetry in the Universe and the mechanism generating the neutrino masses are among the great unsolved questions of modern particle physics. Neutrinoless second-order weak decays are one of the experimental channels available to address these questions by testing the Majorana nature of neutrinos [1][2][3][4]. Most experimental effort to date has focused on searching for the neutrinoless double beta decay (0νβ − β − ) of neutron-rich candidate isotopes [5][6][7][8][9], due to their relatively high natural abundance compared to proton-rich candidates. However, proton-rich isotopes offer unique decay topologies that make them of considerable experimental interest as well. In particular, those with Q-values greater than 2044 keV (4m e c 2 ) can decay in three possible modes -double-electron capture (0/2νECEC), double-positron emission (0/2νβ + β + ) and single electron-capture with coincident positron emission (0/2νECβ + ) [10] -which each produce a different experimental signature. In detectors with high-fidelity position reconstruction, tagging the specific combinations of emitted particles would be a powerful tool for discriminating signal events from backgrounds, potentially providing an extremely low-background or background-free *<EMAIL_ADDRESS>experiment.
While searches for the neutrinoless decays can complement 0νβ − β − searches [2,3] positronic, second-order decays with neutrino emission are theoretically wellestablished [11]. Here, new measurements can be used as a benchmark for nuclear matrix element calculations at the long half-life extreme.
The isotope 124 Xe is of particular interest as its Qvalue of (2856.73 ± 0.12) keV [12] energetically allows all three two-neutrino and neutrinoless decay modes. Its double-K-electron capture (2νECEC) has recently been measured with the XENON1T Dark Matter detector [13]. At T 2νKK 1/2 = (1.8 ± 0.5 stat ± 0.1 sys ) × 10 22 yr the measurement agrees well with recent theoretical predictions [14][15][16]. In this decay, the measurable signal is constituted by the atomic deexcitation cascade of X-rays and Auger electrons that occurs when the vacancies of the captured electrons are refilled. In the XENON1T measurement this cascade was resolved as a single signal at 64. 3 keV. An observation of the KL-capture and LL-capture [10] could be within reach in future experiments if background levels can be controlled, which would allow the decoupling of the nuclear matrix element from phase-space factors. Furthermore, the discovery potential for the positronemitting modes (2νECβ + or 2νβ + β + ) in future, longerexposure experiments could be enhanced by their distinct experimental signatures [17]. Position-sensitive detectors could tag the γ-rays emitted by the annihilating positron, providing a tool for rejecting γ-ray and β-decay backgrounds which arise from natural radioactivity. In beyond-the-Standard-Model neutrinoless decays, the entire energy must be emitted in the form of charged particles or photons, favoring the positron-emitting decay channel 0νECβ + [18][19][20][21][22][23]. As in the two-neutrino case, the coincidence signature of the atomic relaxation, the mono-energetic positron and the two subsequent back-toback γ-rays could be used to reject background. 124 Xe may also allow a resonant enhancement in 0νECEC to an excited state of 124 Te [12], which would be needed to provide accessible experimentally half-lives [24]. The experimental signature contains multiple γ-rays emitted in a cascade, so coincidence techniques can be used to increase experimental sensitivity by suppressing the background substantially.
Liquid xenon time projection chambers (TPCs) are ideally suited to search for 124 Xe decays, due to their large relatively target masses with 1 kg of 124 Xe per tonne of natural xenon, low backgrounds, O(1 %) energy resolution at Q = 2.8 MeV, and position reconstruction for individual interactions within an event. In this work, we investigate the detection prospects of 2νECβ + , 2νβ + β + , 0νECEC, 0νECβ + and 0νβ + β + in multi-tonne xenon TPCs such as the next-generation Dark-Matter detectors LZ [25], PandaX-4t [26] and XENONnT [27], as well as the future nEXO [28] double-β decay experiment, and the DARWIN [29] Dark Matter detector. We simulate the experimental signatures of the second-order 124 Xe decays in such detectors, compute the expected signal detection efficiencies, assess background sources, and calculate the experimental sensitivity as a function of the 124 Xe exposure. We close with a brief discussion on the physics case for pursuing these efforts.
The signal modeling and estimated half-lives of 124 Xe are discussed in II. Relevant details of liquid xenon TPCs are described in III. The detection efficiencies for the different decay channels will be affected by a given detector's energy resolution, spatial resolution, energy threshold and exposure. We outline the analysis of these effects and give the resulting efficiencies in section IV. Potential backgrounds and their impact are discussed in section V. The experimental sensitivities are then given in VI and followed by the discussion in section VII.
II. SIGNALS FROM 124 XE DECAY
The decay modes under investigation provide distinct signatures that can be measured by the coincidence and magnitude of energy depositions (Table I) in a detector. We group the decay modes by the number of emitted positrons. Each emitted positron will lead to the emission of at least two γ-rays and reduce the energy that is initially available for the positrons and neutrinos by twice the positron mass. Each of the 0ν decays will exhibit a monoenergetic total energy deposition while the 2ν decays have continuous spectra due to the neutrinos leaving the detector without further interaction.
We only consider decays to the ground state of the daughter nucleus for the positronic decay modes. A special treatment is required for 0νECEC, as only decays which resonantly populate an excited state of 124 Te may be experimentally accessible.
A. Signal models of decay modes The electron capture with coincident positron emission can be written as where the Standard Model decay features the emission of two electron-neutrinos (ν e ) in addition to the positron (e + ). We assume the most-likely case of an electron capture from the K-shell. This will produce a cascade of X-rays and Auger electrons (X k ) with a total energy of (31.8115 ± 0.0012) keV [30]. The total available energy for the e + and the two ν e is then given by where one has a monoenergetic positron for the neutrinoless decay and a β-like spectrum for the two-neutrino decay. Upon thermalization the e + annihilates with an atomic electron resulting in two back-to-back 511 keV γrays 1 .
The reaction equation for the β + β + -decay to the ground state is 124 Xe → 124 Te + 2e + (+2ν e ). ( The energy available for the two e + and the two ν e is given by where one has a continuous spectrum for the energies of the two positrons for the two-neutrino decay and a 1 The electron mass uncertainty of 44 ppb and the uncertainty on the K-shell X-ray energy in xenon are neglected in our calculations, as they will not affect the results. Moreover, we note that the 2γ-annihilation is by far the most likely case for positronium, but more γ-rays are possible.
Resonant 0νECEC
In contrast to the former decay modes, the energy released in the 0νECEC decay has to be transferred to a matching excited nuclear state 124 Te * of the daughter isotope, since no initial quanta are emitted from the nucleus. For a double-K capture one only has the atomic deexcitation cascade (X 2k ): 124 Xe + 2e − → 124 Te * + X 2k , 124 Te * → 124 Te + multiple γ. (5) The corresponding energy match has to be exact within uncertainties to avoid a violation of energy and momentum conservation. Therefore, the excitation energy E exc,res of the state 124 Te * has to fulfill the resonance condition E exc,res = Q − E 2K = (2856.73 ± 0.12) keV − (64.457 ± 0.012) keV = (2792.27 ± 0.13) keV.
Here, E 2k = (64.457 ± 0.012) keV is the energy of the double electron hole after a double-K capture [12] that occurs in 76.6 % of all decays [20]. The resonance is approximately realized with a positive parity nuclear state at E exc,res = (2790.41 ± 0.09) keV and a corresponding deviation of (1.86 ± 0.15) keV 2 [12,31]. The angular momentum of this state is not precisely known, but 0 + to 4 + 2 The authors of [12] recommend to perform at least one more independent measurement of the 124 Xe→ 124 Te Q-value in order to resolve discrepancies between existing measurements. In addition a determination of J P of the (2790.41 ± 0.09) keV excited state would be helpful in order to further assess the feasibility of this decay mode. are possible J P configurations. The level scheme relevant to the decay is shown in shown in Fig. 1. There are five different γ-cascades that are either ≥ 0 + → 2 + → 0 + or ≥ 0 + → 2 + → 2 + → 0 + for two-and three-γ transitions, respectively. As a considerable decay rate is only expected to 0 + and 1 + states [12], we assume that the resonantly populated state is 0 + and focus on the 0 + → 2 + → 2 + → 0 + transition that occurs in 57.42 % of all decays.
B. Half-life calculations
Two-neutrino decays
The half-life predictions for the two-neutrino decay modes are constructed from where G 2ν is a phase-space factor (PSF) and |M 2ν | 2 is the nuclear matrix element (NME). While the PSF is different among the decay modes [10,18,32], the NME differs only slightly between 2νECEC and 2νECβ + and is about a factor of two smaller for 2νβ + β + [14,15]. For simplicity, we assume and use the existing 2νECEC measurement to constrain M 2νECEC . This is justified by the relatively large uncertainty from the measured half-life [13] which outweighs the expected NME differences. As only the value for the double K-capture has been reported, the half-life has to be scaled by the fraction of double-K decays f 2νKK = 0.767 [20]. One obtains a total half-life of T 2νECEC Using Eq. (7) with the measured half-life and calculated PSFs one has The resulting expected half-lives for 2νECβ + and 2νβ + β + are given in Tab. II. Due to the smaller available phase-spaces, the 2νECβ + half-life is about one order of magnitude longer than the one for 2νECEC, whereas the 2νβ + β + half-life is about six orders of magnitude longer. This makes 2νECβ + a promising target for nextgeneration experiments such as LZ or XENONnT while the double-positronic mode will be challenging to measure.
Neutrinoless decays
In case of the neutrinoless decays the equation relating PSF and NME to the half-life changes to Note that the PSF (G 0ν ) and NME (|M 0ν | 2 ) are different from those used previously due to the absence of neutrino emission. The additional factor f (m i , U ei ) contains physics beyond the Standard Model. Typically the decay is assumed to proceed via light neutrino exchange, for which we have Here the effective neutrino mass m ν is a linear combination of neutrino masses m i and elements of the PMNS mixing matrix U ei [24,33]. For 0νECEC a resonance factor R has to be added to Eq. (11): The mismatch ∆ = |Q − E 2k − E exc | = (1.86 ± 0.15) keV between the available energy and the energy level of the daughter nucleus in the excited state E exc [12] defines the resonance factor R, which -with the two-hole width Γ = 0.0198 keV [24] -amounts to We take the PSF values again from the review [32] which summarizes work by the reviewers and from [20,33]. In order to calculate half-life expectations for neutrinoless decays of 124 Xe, we also need estimates for the NME, and the effective neutrino mass m ν . The NMEs have never been measured for the neutrinoless case. Only for the case with two neutrinos a few half-lives have been determined experimentally. Unfortunately the NMEs for the 2ν and the 0ν cases are not strongly connected. Moreover, the effective neutrino mass has never been measured, and we must choose among different experimental constraints accordingly. To account for these two sources of unknowns we use the following two approaches to get lower limits for the expected half-lives of neutrinoless double-weak decays of 124 Xe.
Method 1: In the first approach, to constrain the effective neutrino mass we take the newest result from the neutrino mass experiment KATRIN which set the most stringent direct, model-independent limit on m ν < 1.1 eV (90% C.L.) [34]. We then combine this limit with a global fit to neutrino oscillation results [35] (Fig. 11 therein). This yields an upper limit range -corresponding to the uncertainties in the Majorana and CP-phases of the PMNS neutrino mixing matrix -on the effective neutrino mass: For the NMEs in our first approach, we take three available sets of calculations into account. The first set is based on the quasi-random phase approximation (QRPA) and was calculated in [14]. The second comes from the interacting boson model (IBM) [24,36]. The third set is based on nuclear shell model (NSM) calculations as performed for the two-neutrino case [16] and is limited by lower and upper values of the full shell model similar to normal neutrinoless double-β decay as shown in [37] and [38]. Both the QRPA and NSM calculations provided good predictions of T 1/2 for 2νECEC while there were no 2ν-predictions for IBM. We summarize the relevant PSF-and NME-values and the corresponding lower half-life limits in Tab. III. Method 2: In our second approach, we use a similar idea as for the prediction of the half-lives in the 2ν case. Instead of one measured half-life value we take the halflife limits obtained in the search for 0νβ − β − decay of II: The different 2ν decay modes of 124 Xe with the corresponding phase-space factors (PSF), the assumptions of the corresponding matrix elements according to Eq. (8), and the measured or predicted half-lives according to Eq. (9) and Eq. (10), respectively. The PSF values were taken from the review [32] which summarizes work by the reviewers and from [10,33]. Therefore, we give a range of PSF values. For predicting the half-lives of the decay modes 2νECβ + and 2νβ + β + we use the central value of this range as the most probable PSF value and half of this range as the uncertainty.
136 Xe [5,6]. The most stringent lower limit on the halflife comes from the KamLAND-Zen experiment [5] with Unlike the case for the various 2ν-decays in Eq. (10), the NMEs of 0ν-decays of 124 Xe are different from the NMEs of the 0νβ − β − -decay of 136 Xe and do not cancel. But for this comparison of the half-life limits of neutron-poor and neutron-rich isotopes of the same element xenon, the substantial uncertainties connected to these calculations drop out to a large extent if the NMEs are calculated within the same framework and if the main uncertainties stem from the unknown quenching q of the axial coupling constant g A , which can be factorized out of the NME M : This is of advantage, as the uncertainties of the NME are often summarized in the g A -quenching, and the quenching factors q can be assumed to be similar for neutronpoor and neutron-rich nuclei of the same element [39]. We perform the calculations using the NMEs from the interacting boson model (IBM) [24,36] which possess their main uncertainties in the quenching of the axial coupling constant. If we solve equation (11) or (13) according to the factor f (m i , U ei ) and equate for two decays of different xenon istopes, we obtain: T 0νECEC,124 In Eq. (18) the effective neutrino mass and its uncertainty drop out of the equation. We would like to underline the validity of this kind of comparison by mentioning that a similar approach has been used in Ref. [40] (Fig. 6 therein) to relate effective neutrino mass limits for two different isotopes within the same theoretical framework of NME calculations. Again, we summarize the results in Tab. IV. Among the neutrinoless decays, the 0νECβ + mode is expected to have the shortest -and thus most experimentally accessible -half-life. The other decay modes, 0νβ + β + and resonant 0νECEC, exhibit considerably longer half-lives owed to unfavourable phase-space and a lack of resonance enhancement R, respectively. We also note that the half-life limits calculated with the first method are systematically lower than the ones for the second. This is in line with the smaller predicted upper limits on the effective neutrino mass given by KamLAND Zen [5]. The ability to observe any of the described decay channels is given not only by the theoretical prediction on their half-lives, but also by the detection efficiencies in a given experiment. In the following sections, we discuss the detection prospects of the neutrinoless decay modes in future experiments which could have significantly larger samples of 124 Xe than current detectors, such as DARWIN or nEXO.
III. DETECTOR RESPONSE
In a liquid xenon TPC, the energy and position of an energy deposit is reconstructed using two observed signals: scintillation light and ionization charge [43]. The former is typically detected directly using UV-sensitive photodetectors, producing a prompt signal referred to as S1. The latter is detected by applying an electric field across the liquid xenon volume and drifting the charges to a collection plane. The charge can be either detected directly using charge-sensitive amplifiers, or extracted into a gas-phase region and accelerated, producing proportional electroluminescence light that is detected in the [34,35]. The PSFs (G 0ν ) were taken from [24], and the review [32] which summarizes work by the reviewers and from [18,20]. We use the central value of the PSF-range as the most probable value and half of this range as the uncertainty. The same is done for the NMEs (M 0ν ) in all cases were a range of values is given in the original publication. For 0νECEC the NMEs values from the quasi-random phase approximation (QRPA) [14] and the interacting boson model (IBM) [24] were used. The NME for IBM is obtained by taking the single value given in the publication and assuming g A = 1.269. The NME-range for QRPA stems from the smallest and largest NME value for g A = 1.25 under the assumption of different bases and short-range correlations. For the 0νECβ + and the 0νβ + β + QRPA [14], NSM (calculated [41] as in [16]), and IBM [36] NMEs were considered. The range of NMEs for QRPA and the value for IBM are obtained as above. However, for the latter an uncertainty is given in the publication instead of a value range. For the NSM the NME-range is given by different model configurations and the most probable value and uncertainty are derived in the same fashion as for QRPA. All uncertainties are propagated by drawing 10 6 independent samples from the parameter distributions and multiplying with the upper limit on m ν . Then the 90 % C.L. upper limit on T −1 1/2 is determined from the resulting distribution and inverted to obtain the corresponding lower half-life limit. photodetectors. The delayed secondary signal produced by the drifted charge is referred to as S2. The combination of the two signals allows one to reconstruct the 3D-position of the interaction inside the detector: the S2 hit pattern on the collection plane gives the x-y coordinate and the S1-S2 time delay gives the depth z. The deposited energy is reconstructed using the magnitude of the S1 and S2 signals. A linear combination of both signals has been shown to greatly improve the energy resolution compared to either signal individually, due to recombination of electron-ion pairs producing anticorrelated fluctuations in the energy partitioning between light and charge [44,45].
For events with multiple energy deposits, the prompt S1 signals for each vertex are typically merged, resulting in a single scintillation pulse for the entire event. However, individual vertices can be resolved as individual S2 signals arriving at different positions and times on the charge collection plane. A schematic of the signature expected from a typical 0/2νECβ + decay of 124 Xe is shown in Fig. 2. In this example, there are five different S2 signals produced from the positron, X-ray cascade, and each of the annihilation γ-rays (one of which undergoes Compton scattering before being absorbed). With sufficient position and energy resolution, one can use this information to classify events and perform particle identification, providing a tool for separating backgrounds from the signal of interest.
The capability for a detector to resolve each vertex depends on the time resolution in the charge channel, the width of the S2 signals, and the x-y resolution of the charge collection plane. These properties are highly dependent on the specific readout techniques employed in each experiment. In addition, the detection of each energy deposit requires its individual S2 signal to be above the detector's charge energy threshold, a property that is again specific to each experiment.
In this work, we compute the detection efficiency ( ) for the various modes of 124 Xe decay as a function of the x-y-and z-position resolution and energy threshold, to provide estimates that apply across the possible range of existing and future experiments.
A. Simulation
We generate the emitted quanta and their initial momentum vectors for each decay channel with the event ) were taken from [24], and the review [32] which summarizes work by the reviewers and from [18,20]. Those for 136 Xe (G 136 Xe 0ν ) were also taken from [32] and [42], cited therein. We use the central value of the PSF-range as the most probable value and half of this range as the uncertainty. . For 0νECEC the NME was taken from [24]. For 0νECβ + and 0νβ + β + the NMEs were taken from [36]. An uncertainty on the NME is only given in this publication. All uncertainties are propagated by drawing 10 6 independent samples from the parameter distributions and multiplying with the lower limit on T 0νβ − β − ,136 . Then the 90 % C.L. lower limit of the half-life is determined from the resulting distribution. generator DECAY0 [46]. The version used here has been modified previously for the simulation of the positronic 124 Xe decay modes [17]. In the scope of this work, we verified the implementation, added the resonant 0νECEC decay mode, and implemented the angular correlations for the γ-cascades under the assumption of J P = 0 + for the resonantly populated state [47,48]. In order to investigate the efficiency, at least 10 4 events per decay channel have been used. The particles generated for each decay are propagated through simplified models of the detectors under investigation using the XeSim package [49], based on Geant4 [50]. These detector models consist of a cylindrical liquid xenon volume in which we uniformly generate 124 Xe decay events. This volume is surrounded by a thin shell of copper which is used for modeling the impact of external γ-backgrounds. We simulate two different sizes of cylinders in this work, characteristic of two classes of future experiments. The "Generation 2" (G2) experiments are defined as experiments which have height/diameter dimensions of between one and two meters. This class includes the LZ [25] and XENONnT Dark Matter experiments, which will use dual-phase TPCs filled with natural xenon. It also includes the future nEXO neutrinoless double-β decay experiment, which will use a singlephase liquid TPC filled with xenon enriched to 90% in 136 Xe. For simplicity, we model all G2 experiments as a right-cylinder of liquid xenon with a height and diameter of 120 cm each. 3 We also simulate a "Generation 3" (G3) experiment, which is intended to model the proposed DARWIN Dark Matter experiment [29]. This de-tector is modeled as a right-cylinder of liquid xenon with a height and diameter of 250 cm each.
For experiments using nat Xe targets, there will be approximately 1 kg of 124 Xe per tonne of target material. The G2 Dark Matter experiments would therefore be able to reach 124 Xe-exposures of ∼ 50−100 kg-year in 10 years of run time. By scaling the target mass up to 50 tonnes, the G3 experiment DARWIN will amass an exposure of ∼500 kg-year. For nEXO, the enrichment of the target in 136 Xe will remove all of the 124 Xe; however, here we consider the possibility of extracting the 124 Xe from the depleted xenon and mixing it back into the target. There will be approximately 50 kg of 124 Xe in the nEXO tailings, meaning a 10 year experiment could amass an exposure of ∼500 kg-year, competitive with a G3 natural xenon experiment.
B. Energy resolution model
Within this study all simulated detectors use the energy dependence of the resolution on the combined signal as reported in [13], which is modeled as Here E is the energy and a = 31 keV 1 /2 and b = 0.37 are constants extracted from a fit to calibration data from 41−511 keV. The model predicts a resolution of ∼ 1 % at the Q-value of the decay, approximately consistent with the energy resolution published by the EXO-200 experiment at a similar energy [51]. This energy resolution is used to define the full-energy region of interest (ROI) for the various modes of 124 Xe decay. For the neutrinoless modes, this corresponds to a narrow energy window around the Q-value. For the filtering of single energy depositions within an event, the reconstruction can only be based on the S2. To model the broadening of the charge-only energy resolution due to recombination fluctuations, we scale b in the above formula to a value of 4.4. This gives a chargeonly resolution of about 6 % at ∼500 keV, consistent with measurements reported in the literature [52,53].
C. Event reconstruction and efficiency calculation
This analysis utilizes the information on the various energy depositions at a given spatial position in all three dimensions. In order to reconstruct and validate the efficiency for detecting the unique event topologies, several filtering and clustering steps have to be performed 4 .
First the events are filtered by the total energy deposited in the detector, in order to account for events where decay products leave the detector. This criterion is a fixed value, only broadened by the energy resolution for the neutrinoless modes, but a broad range with a maximum cut-off at the Q-value and a decay-dependent threshold for the two-neutrino decays.
For any remaining event the vertices are sorted by their axial position in the detector (z-coordinate) and these vertices are grouped within a spatial range determined by the assumed position resolution of the detector in the axial direction. For detector configurations where a separation in the radial direction (x-y-coordinate) is also possible, the grouping algorithm also takes separations in x-y into account -according to the assumed position resolution. The energies of all vertices within each group are summed and provide the individual S2 signals that a detector with the chosen properties would see.
From this point the further filtering targets the reconstruction of the vertices of the annihilation products of the positrons. The procedure is analogously applied to the de-excitation γs in the case of the 0νECEC. It is depicted together with an illustration of the spatial clustering in Figure 2. All clustered energy depositions of a given event are permuted for each possible interaction combination and the total sum of the energy is compared against the expected value, which is e.g. 511 keV for each γ produced in the positron's annihilation. The combination with the smallest difference between the summed energy and the expected value is then removed from the list of energy depositions if it lies within the energy resolution around the expected value. This raises the counter of measurable signatures by one. Afterwards, this procedure is repeated until all desired signatures have been found and the counter matches the expectation (e.g. 4 in 2νβ + β + ). For any left-over energy it is then checked if it fulfills the requirement for the point-like deposition expected from the positron and/or the electron capture signal. In case of 0νECβ + /2νECβ + a single merged energy deposition of the positron and atomic relaxation processes is expected. While this requirement is a fixed maximum value for a single signature in case of the neutrinoless mode, it is again a continuous distribution ranging from zero or the 31.81 keV K-shell hole energy to a The efficiency for the 0νECβ + (blue) and 2νECβ + (black) show a decrease with energy up to about 250 keV with efficiencies ranging from about 41 % to 10 %. More striking is the behavior of the 0νECEC (red), which has a sharp cut-off as soon as the double electron capture energy (64.3 keV) is below the threshold. Since this signature is required within this analysis in order to provide a clear evidence and necessary background suppression, this will automatically drop the efficiency to zero. For this example a position resolution of 10 mm has been used in both directions.
cut-off depending on the Q-value. The requirement removes energy signatures which are merged by the detector due to the aforementioned limited spatial and time resolution. If not all signatures have been found or if the remaining energy is not a single deposition, the event is discarded. The ratio of all events which survive the filtering algorithm and the original generated number of events corresponds to the desired efficiency .
D. Influence of thresholds, detector size and position resolution
We first investigate the effect of the S2 energy threshold on the detection efficiency. While for Dark Matter detectors the threshold usually is only a few keV thanks to amplification via electroluminescence, the situation is different for an experiment like nEXO, which will measure charge directly. In this case, the electronics noise in the readout circuit introduces a larger energy threshold and thus influences the efficiency for the example decay modes as shown in Figure 3. In this work, the threshold is implemented in the simulation simplified as a sharp cutoff for any given energy signature, assuming a position resolution of 1 cm in radial and axial direction. It is evident that the efficiency depends on this energy threshold. Therefore, an improvement from O(100 keV) as achieved in EXO-200 [54] would be beneficial for nEXO as this has a direct impact on the sensitivity for any given decay channel. Especially, in order to look for a possible smoking gun evidence of the 0νECEC decay, a threshold below the energy of twice the K-shell electron energy is necessary. In the following we assume that a sufficiently low threshold is achieved that it can be considered negligible.
Next, we investigate the effect of a detector's position resolution on the detection efficiency. Our results are shown in Figure 4 where we emphasize the importance of x-y resolution. For any detector with an axial position resolution (z-coordinate) of a few mm, which is fundamentally limited by electron diffusion, an additional resolution of event topologies in the radial direction is highly beneficial. Already at an achieved 10 mm separation in the axial direction, an x-y resolution of also 10 mm can improve the efficiency by more than a factor of two. For a nEXO-type detector this resolution is mostly a function of the pitch of the charge readout strips [55], and therefore can become as small as a few mm. The situation is less clear for dual-phase detectors used in Dark Matter searches; no detector dedicated for Dark Matter search has reported its x-y resolution for multiple energy depositions arriving at the charge detection plane simultaneously. In principle this should be achievable by pattern recognition in the top array of the detector, and is a good candidate for future work in better matching algorithms and machine learning techniques.
Finally, an interesting comparison arises between a nEXO like detector and a G3 Dark Matter experiment, as both could have the same amount of 124 Xe within different-sized detector volumes. The influence of the detector size on the efficiency for the decay mode of 0νECβ + is shown in Fig. 4. It is evident that an increased detector size only increases the efficiency by a few %. This is due to the ratio of events leaving the detector in comparison to the events confined in the full volume. Therefore, the findings for a G3 detector that are summarized in Table V are approximately also valid for a nEXO-like detector.
V. BACKGROUNDS
From the above analysis, it is clear that the most experimentally accessible decay channels are the 0ν/2νECβ + . As described, the key feature in a search for β + -emitting decay modes is the ability to reject backgrounds using the distinct event topology. We consider possible sources of backgrounds below and estimate the expected rates of events passing the topological selection criteria described in Section IV.
As comparison points, we compute the expected number of 124 Xe decays per tonne-year exposure of nat Xe (corresponding to 0.95 kg-year of 124 Xe) using the halflives estimated in Table II, Table III and Table IV. After including the respective efficiencies for a G2 exper- a Here we considered the most probable branch (57.42%) with a three-fold γ-signature. An analysis using the two-fold signatures would yield higher efficiency but can add coincidental γ-backgrounds, which would weaken the sensitivity of a given search. Efficiency iment with 10 mm resolution in x-y-z and assuming a natural xenon ( nat Xe) target, we expect 8.3 ± 2.9 decays per tonne-year for 2νECβ + . Under the assumption of light-neutrino exchange and given the most optimistic assumptions described above, we expect a rate of less than 2.6 · 10 −2 decays per tonne-year for 0νECβ + .
A. Radiogenic backgrounds from detector materials
Gamma rays from radioactivity in the laboratory environment and detector construction materials are a primary background in rare event searches. There are two main concerns for the analysis presented here: first, that a γ-ray Compton-scatters multiple times and produces the expected event signature. Second, that a γ-ray of sufficient energy creates a positron by pair production. In the latter case, the positron will annihilate and produce a background event which, by design, passes our event topology cuts.
We investigate the sources for falsely identified events from the 238 U and 232 Th decay chains, the most common sources of radiogenic backgrounds in most 0νββ searches. For each decay step within the chain, 10 7 events 5 have been uniformly generated in a copper shell of 1 cm thickness surrounding the liquid xenon volume of a G2-sized detector using Geant4. Afterwards, the events which interacted in the active volume were run through the respective event search algorithms for 2νECβ + and 0νECβ + . We find that the only relevant decays are βdecays into excited daughters, as only these produce γs of sufficient energies.
For the neutrinoless case there are two particularly problematic transitions. The first is the β-decay of 214 Bi in the 238 U-chain, which has a small branching to the 2880 keV state of 214 Po. If this γ-ray interacts via pair production, it creates an event identical to our signal directly in the ROI. We find that 1.5 · 10 −6 events per 214 Bi primary decay pass the selection criteria. The second problematic transition is the decay of 208 Tl to 208 Pb in the 232 Th-chain, for which there are various transitions in which different γ-rays are detected in coincidence with the one from the 2614 keV state. Such events can deposit enough energy to create events in the ROI, and may similarly produce a sequence of energy depositions which pass our topological criteria. We find that 1.3 · 10 −4 events pass our cuts per 208 Tl primary, but the 35.9 % branching fraction for creating 208 Tl in the first place reduces its impact in a real detector to 4.5 · 10 −5 events per 232 Th primary decay. Both sources of background can be reduced by a subselection of an inner volume in the active volume of the detector, commonly referred to as a "fiducial volume cut." As different γ-rays with energies below 300 keV are paired with a high energy γ in the 208 Tl decay signature, fiducializing is especially effective against these events; we find that cutting away the outer 10 cm of LXe reduces its background contribution by almost an order of magnitude. For a 20 cm cut, no event out of the 10 7 simulated for any isotope passes the selection criteria. We conclude that these backgrounds can therefore be eliminated in a real experiment (depending on the actual 238 U/ 232 Th contamination) by selecting an appropriate fiducial volume.
Radiogenic backgrounds have a greater impact on 2νECβ + searches, as the larger energy window allows more events to pass the selection criteria. We find three isotopes in the 238 U chain producing events which pass our selection criteria, with decays of 214 Bi into different excited states of 214 Po being the major background component (>99 %). The surviving fraction for the total chain is 6.9 · 10 −3 events per 238 U primary decay without a fiducial volume selection. This is reduced to 1.5 · 10 −3 and 1.1 · 10 −4 decays per primary with the 10 cm and 20 cm cuts respectively. For the 232 Th-chain, the 208 Pb γ-rays following 208 Tl β-decay are again the main contributor (∼ 75%). However, γ-rays from 228 Th after the β-decay of 228 Ac also contribute (∼ 23%), as well as a small contribution (∼ 2%) from excited states of 212 Pb following the β-branch of the 212 Bi decay. The surviving fractions for the whole chain are 7.3 · 10 −3 , 1.5 · 10 −3 and 1.3 · 10 −4 events per primary 232 Th decay with no fiducial volume cut, a 10 cm cut and a 20 cm cut, respectively. Due to the less-stringent energy selection the fiducial volume cuts are less efficient for the 208 Tl events in the two-neutrino case, but still significantly reduce the background contribution.
In conclusion, two factors play a role for the exact evaluation of a given experimental setting: the fiducial volume cut and the actual amount of contaminants surrounding the TPC. While this study cannot provide an answer for all given experimental settings -this would need a dedicated Monte Carlo study following a material radioassay -we use reported contamination levels and experimental details projected for the nEXO experiment (reported in Ref. [56]) to benchmark our calculations. Our approximate evaluation of a nEXO-like experiment is provided in Fig. 5. The nEXO experiment identifies the main source of external γ-ray backgrounds as the copper cryostat, for which the collaboration reports 238 U and 232 Th concentrations of 0.26 ppt and 0.13 ppt, respectively. This corresponds to 2.8 · 10 5 primary decays per year as indicated in Figure 5 by the dotted gray line. Accordingly, it would only require a mild 10−20 cm fiducial volume cut to achieve a favorable signal to background ratio 6 hand, are optimized for the low-energy regime, and typically have higher background levels in the ∼MeV regime, and may therefore require more aggressive fiducial cuts to achieve a similar signal-to-background ratio.
We emphasize that these results are only approximate, and that in a full likelihood analysis the modeling of the events' spatial components and energy distributions will improve the signal to background ratio beyond what has been discussed above. More precise estimates of the impact of these backgrounds are beyond the scope of this work, but will be necessary to understand the true impact of externally-produced γ-ray backgrounds in real experiments.
B. 222 Rn 222 Rn may dissolve into the active LXe volume and create backgrounds via β-decays that emit γ-rays (α-decay events can be easily rejected by the ratio of ionized charge to scintillation light [57]). There are only two β-decays in the 222 Rn chain with enough energy to create backgrounds in this analysis. The first, 214 Bi, is accompanied by the subsequent α decay of 214 Po, which occurs with T 1/2 = 164 µs. Thus, we assume that it can be rejected via a coincidence analysis. The second, 210 Bi, has a Q-value of 1.2 MeV -just at the low-energy end of our region of interest for the 2ν decays, but well below the ROI for 0ν signals -and decays with no accompanying γ. Therefore, it almost always is a single-scatter signal and does not pass our cuts.
C. Charged-current scattering of (anti)neutrinos Charged-current (CC) scattering of neutrinos and antineutrinos, while rare, may produce positrons which can exactly mimic our signal of interest.
The CC scattering of low-energy antineutrinos produces a fast positron in the final state. Here we consider two sources of antineutrinos: nuclear reactors and radioactive decay in the earth (geoneutrinos). Both of these are sources of electron antineutrinos in the few-MeV range. The threshold for the charged-current reaction is set by the mass difference between the xenon isotopes and their iodine isobars. The cross-sections as a function of energy were computed in Ref. [58], and were obtained in tabular form from the authors.
We calculate the expected rates for geoneutrinos using the two xenon isotopes with the lowest CC reaction thresholds: 129 Xe and 131 Xe, which have thresholds of 1.2 MeV and 2.0 MeV, respectively. Convolving the energy spectra and flux with the cross section, we find that the rates for 129 Xe and 131 Xe are 5.0×10 −8 and 4.9×10 −6 events per tonne-year of nat Xe exposure, respectively. In a G3 detector filled with nat Xe, there will, therefore, be less than 0.01 events in a 10-year exposure, rendering this background negligible. An experiment using xenon enriched in the heaviest isotopes ( 134 Xe and 136 Xe) will be completely insensitive to geoneutrinos due to the high threshold for CC reactions. The flux of geoneutrinos is expected to vary by a factor of ∼2 across the globe, so we do not expect these conclusions to depend on the location of the experiment.
We carry out a similar calculation for reactor antineutrinos, which in contrast are highly location dependent. We assume three possible locations for an experiment: SNOLAB (in Sudbury, Ontario, CA), Sanford Underground Research Facility (in Lead, South Dakota, USA), and Laboratori Nazionali del Gran Sasso (in L'Aquila, Abruzzo, Italy). The reactor antineutrino flux at each site is calculated using reactor power and location data in the Antineutrino Global Map reactor database [59]. The antineutrino flux and energy spectra are computed using the empirical models given in Ref. [60]. For simplicity, we neglect neutrino oscillations, meaning our expected rates will be overestimated. Of the three candidate locations, the flux is highest at SNOLAB, primarily due to the presence of nearby reactors in Kincardine and Pickering, ON. In this case, we calculate an expected CC scattering rate of 9.1×10 −7 and 3.6×10 −6 events per tonne-year for scattering on 129 Xe and 131 Xe, respectively. The expected rates at the other candidate locations are smaller by at least an order of magnitude.
In contrast to antineutrinos, CC scattering of lowenergy neutrinos does not directly create positrons.
There are, however, two possible backgrounds that may arise from this reaction: the emission of a fast electron and a daughter nucleus in an excited state (which can de-excite and create additional energy deposits that may mimic the signal event topology), and the creation of a daughter radioisotope which later decays via β +emission.
The primary source of neutrinos incident on deep underground detectors is the sun. For our purposes, the most important are those produced by the decay of 8 B, which have energies of ∼ 1 − 10 MeV. These are the only solar neutrinos with enough energy to react above threshold and populate an excited state in the daughter nucleus, for all xenon isotopes. The energy-averaged crosssection for these reactions is tabulated in Ref. [58], and is of O 10 −42 − 10 −41 cm 2 . This may produce 10's of events per tonne-year for each isotope in a nat Xe detector. However, scattering into low-lying excited states is suppressed, with partial cross-sections an order of magnitude smaller than the total. Therefore, most of the events will deposit too much energy in the detector and will be rejected. We find that the low probability of the remaining events passing our topological selection criteria renders these backgrounds negligible for both the 2νECβ + and the 0νECβ + decay modes.
Neutrino CC scattering on xenon may create radioactive isotopes of caesium in the liquid target. Of particular concern are 128 Cs and 130 Cs, which each have half-lives of < 1 hr and can decay via β + -emission with Q-values of 3.9 MeV and 2.9 MeV, respectively, exactly mimicking our expected event signature. Again using the 8 Baveraged cross sections from Ref. [58], we calculate a production rate of 0.02 nuclei of 128 Cs and 0.07 nuclei of 130 Cs per tonne-year of nat Xe exposure. The resulting β + decays are distributed across a broad spectrum, and our simulations indicate that they will be a small background for the 2νECβ + process, with expected rates an order of magnitude lower than the expected signal rate. The narrow ROI for 0ν searches will render these backgrounds negligible. There are also two isotopes of xenon with CC reaction thresholds low enough to react with CNO, 7 Be, and pp neutrinos: 131 Xe and 136 Xe. However, the relevant Cs daughter isotopes have half-lives of O(10) days. Next-generation experiments plan to recirculate and purify the liquid xenon with a turnover time of ∼2 days [55], meaning these isotopes will likely be removed from the detector before they decay.
D. Neutron-induced backgrounds
A final possibility for backgrounds are those from neutron scattering or capture. In neutron capture, the daughter nucleus is generally left in a highly-excited state, and relaxes to the ground state via the emission of several γ-rays. As the sum total of the energy lost in this process is well above the Q-value for 124 Xe decay, we expect these events will be easy to reject and we neglect this as a background source.
For neutron scattering, it is of particular interest to estimate the activation rate of Xe-radioisotopes that may decay via β + emission in the region of interest. We identify the fast neutron scattering 124 Xe(n, 2n) 123 Xe reaction as the only one of significance. It has a neutron-energy threshold of 10.5 MeV and the crosssection reaches ∼1.4 barn at a neutron energy of ∼ 20 MeV [61]. The high threshold prevents radiogenic neutrons (which come from (α,n) reactions in the laboratory environment) from producing this background, but muon-induced neutrons, which can extend in energy up to the GeV scale, are of concern. We use an estimate of the muon-induced neutron flux at Gran Sasso of 10 −9 n/cm 2 /s and multiply by a factor of 10 −2 to account for the expected reduction due to shielding typically employed in these experiments [62]. We find an expected activation rate of ∼10 −3 atoms per kg ( 124 Xe) per year, each of which we assume will produce a background event in the TPC. However, this decay has a small branching ratio for β + decay, and even then always proceeds to an excited state of 123 I. Accordingly, it is unlikely to pass our selection criteria, and we consider it negligible.
E. Summary
After considering an exhaustive list of background sources, for 2νECβ + we conclude that the only significant background originates from external γ-rays. With strong fiducial volume cuts, a likelihood-analysis utilizing energy information and γ-background suppression, near-future G2 Dark Matter experiments have a strong chance of measuring this decay mode. For 0νECβ + , we conclude that the searches in G2 and G3 experiments will basically be "background-free," and the sensitivity will only be limited by the detection efficiencies and the attainable 124 Xe exposure in each experiment.
VI. SENSITIVITY
The half-life measured by a detector configuration with no expected background for a number of N observed decay events is given by Here, N A is Avogadro's constant, is the detection efficiency, and M Xe corresponds to the molar mass of 124 Xe. The available mass of 124 Xe, m, and the measurement time, t, depend on the detector configuration. If no events are observed and if a Poissonian process without background is assumed, a 90 % C.L. lower limit on T 1/2 can be calculated by inserting N = 2.3. For a detector with 10 mm position resolution in the axial as well as the x-y direction, the expected half-life can be calculated as a function of exposure using the previously calculated efficiencies. The sensitivities for a G3 experiment with a 500 kg-year exposure are summarized in Table VI for all decay modes. A similar exposure would be possible in a G2 detector enriched to 50 kg of 124 Xe; the only difference is the ∼ 10% decrease in detection efficiency due to the increased probability of energy being deposited outside the sensitive volume of the detector. The sensitivities are compared to the range of theoretical predictions from Table II for 2ν-decays, and Tables III and IV for 0ν-decays. Regarding the two-neutrino decays 2νECβ + will likely be detected by a G3 experiment, but are already be accessible to a G2 detector with a nat Xe target if the γbackground is properly addressed. However, due to an unfavourable phase-space 2νβ + β + will likely be out of reach of even a G3 detector. On the neutrinoless side 0νβ + β + is also pushed to experimentally inaccessible half-lives by the unfavourable phase-space. An eventual detection of 0νECEC relies on the presence of a sufficient resonance enhancement that could boost the decay rate approximately four orders of magnitude. However, given current measurements of decay energies and 124 Te energy levels this is not present [12,31]. A final independent measurement is recommended by the authors of [12] would be needed for a final verdict on the detection prospects of this decay. Thus, it is evident that the most promising neutrinoless decay is 0νECβ + .
For this decay we compare the experimental sensitivity derived in this study with three possible theoretical scenarios. Scenarios one and two are based on the direct calculation (Method 1, Eq. 11 and Eq. 13) using the effective neutrino mass range from Eq. 15. Scenario three is based on the comparison of 124 Xe-NMEs with the NME for 0νβ − β − of 136 Xe using the KamLAND Zen half-life limit (Method 2, Eq. 18). The results are shown as a function of exposure in Fig. 6. Within a 500 kg-yearexposure, a background-free experiment would cover a significant portion of the parameter space given by the KATRIN limit translated to m ν . Once this value is reduced, e.g. by phase cancellations in the PMNS-matrix, the lower limits on the half-life are an order of magnitude above the experimental sensitivity. Assuming the same decay mechanism for 136 Xe and 124 Xe -here lightneutrino exchange -the expected half-lives are two orders of magnitude above the experimental sensitivity taking into account the current limits placed by KamLAND Zen. Exposures larger than 10 4 kg-year would be needed to probe this parameter space.
VII. DISCUSSION
This work has summarized the possible decay modes of 124 Xe and investigated possible efficiencies of future liquid xenon detectors to the respective channels. For a G2 Dark Matter detector a detection of 2νECβ + is feasible given a proper treatment of potential γ-backgrounds. An TABLE VI: Results and theoretical predictions for the various decay channels of 124 Xe. The experimental sensitivity is calculated for a 500 kg-year exposure assuming a G3-experiment with 10 mm position resolution in all three dimensions, a negligible threshold, and no backgrounds. The range of theoretical predictions for neutrinoless decays is given between the weakest limit from the direct calculation with m ν < 1.1 eV (Table III) and the strongest limit from the comparison with KamLAND Zen (Table IV) for a background-free experiment with 10 mm resolution in x-y-z, as a function of the exposure (red). This calculation assumes the G3 geometry; the sensitivity curve decreases by ∼10% for a G2-sized detector at all exposures. Three ranges of lower limits on the 0νECβ + -decay half-life are shown: the direct calculation (Table III) with m ν < 1.1, eV (light blue) and with m ν < 0.3 eV (medium blue), and the NME comparison (Table IV) using the KamLAND Zen 136 Xe 0νβ − β − half-life limit (dark blue). For the direct method the lower bound is given by the weakest limit among the three NMEs for each m ν . experiment with the expected background of a doubleβ decay detector like nEXO, would be able to clearly detect the decay and could study it with precision, if 124 Xe would be added to the xenon inventory. A G3 Dark Matter experiment like DARWIN would have the signal strength to detect this decay with a few thousand signals, but would need to optimize its fiducial volume in order to reduce the γ-background.
For a possible neutrinoless mode of this decay, achieving a background-free experiment is a realistic prospect owed to the decay signatures. However, we have shown that in this case a detection is only within reach of a G3 or an enriched nEXO-like detector for the most conservative half-life predictions. It has to be emphasized that such a scenario would require a mechanism that leads to a difference in the decay of proton-rich nuclei compared to their neutron-rich counterparts. Otherwise it would be excluded by KamLAND Zen.
As mentioned previously, such a mechanism would be an exciting prospect in searches for neutrinoless decays of proton-rich isotopes. If detected, it would provide complementary information on the physical mechanism mediating the decay process. One example for this possibility was studied in detail in Ref. [21] in the context of leftright symmetric models, in which one assumes that there is a right-handed weak sector in addition to left-handed neutrinos, which can mediate neutrinoless double-β decays. Detectors with the capability of measuring both isotopes simultaneously may therefore be attractive for both the discovery of the neutrinoless process and subsequent study of the underlying physics. Here we briefly reexamine the analysis of left-right symmetric models using the projected sensitivities described in this work. By adding right-handed terms to the Standard Model Lagrangian, one derives a new expression for the half-life of neutrinoless second-order weak decays: where α represents the decay mode (0νβ − β − , 0νECβ + , etc.), m ν is the effective light neutrino mass defined above, and η and λ are the effective coupling parameters for the new interaction terms containing righthanded currents. The coefficients C α ij are combinations of nuclear matrix elements and phase space factors, and differ between the decay modes. In particular, it was Here we assume η = 0. The exclusion limits compare the present limits on the 0νβ − β − -decay of 136 Xe [5] with the possible limits on 0νECβ + derived in this work. We assume the full 500 kg-year exposure for the 124 Xe search -comparable to the 504 kg-year exposure used for the 136 Xe measurements. The dashed line represents the boundary of the excluded zone after arbitrarily scaling the NMEs for 124 Xe by a factor of two, to mimic uncertainties in NME calculations.
pointed out in Ref. [21] that the λ terms are significantly enhanced in the case of the mixed-mode decays, meaning the shape of the parameter space explored by 0νECβ + searches differs from that explored by the more common 0νβ − β − experiments. We illustrate this in Figure 7, where we compare the possible limits for 0νECβ + derived in this work with the current limits for the 0νβ − β − of 136 Xe decay from the Kamland-Zen experiment [5]. We see that the sensitivity of the mixed-mode 124 Xe decay to the effective neutrino mass is significantly weaker; this is due to the reduced phase space in the positron-emitting decay mode. However, the sensitivity of the mixed-mode decay is within a factor of two for the right-handed coupling λ , which is within the uncertainties typically assumed for nuclear matrix element calcula-tions (usually a factor of ∼ 3). Consequently, such a measurement would provide complementary information in the event of a discovery of a 0ν decay mode in either isotope. It must be acknowledged that future experiments expect to reach sensitivities considerably larger than the existing limits. Unless the 0νβ − β − decay of 136 Xe is just beyond the reach of present experiments, we show that the 124 Xe mixed-mode decays will not be competitive in constraining left-right symmetric models with a G3 experiment's exposure. However, exploring proton-rich isotopes may still provide complementary information in determining the mechanism of lepton number violation; for example, an (unexpected) discovery of neutrinoless decays in either only 124 Xe or in both 124 Xe and 136 Xe could prove that neither the light neutrino exchange nor right-handed currents mediate the decay processes, and could point towards alternative new physics. Therefore, we emphasize that future xenon-based TPC experiments should explore this decay channel, as the striking multiple coincidence structure is straightforward to look for and distinguish from backgrounds. Also the consideration to expand an existing program like nEXO, which would require an additional enrichment on the light mass side after the initial enrichment, could be thought of, in order to gain further insight into the neutrinoless decay modes -especially once it has been found in 136 Xe. | 13,911.8 | 2020-02-11T00:00:00.000 | [
"Physics"
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Antihyperglycemic Activity of Houttuynia cordata Thunb. in Streptozotocin-Induced Diabetic Rats
Present study is an attempt to investigate plausible mechanism involved behind antidiabetic activity of standardized Houttuynia cordata Thunb. extract in streptozotocin-induced diabetic rats. The plant is used as a medicinal salad for lowering blood sugar level in North-Eastern parts of India. Oral administration of extract at 200 and 400 mg/kg dose level daily for 21 days showed a significant (P < 0.05) decrease in fasting plasma glucose and also elevated insulin level in streptozotocin-induced diabetic rats. It also significantly reversed all the alterations in biochemical parameters, that is, total lipid profile, blood urea, creatinine, protein, and antioxidant enzymes in liver, pancreas, and adipose tissue of diabetic rats. Furthermore, we have demonstrated that the extract significantly reversed the expression patterns of various glucose homeostatic enzyme genes like GLUT-2, GLUT-4, and caspase-3 levels but did not show any significant effect on PPAR-γ protein expressions. Additionally, the extract positively regulated mitochondrial membrane potential and succinate dehydrogenase (SDH) activity in diabetic rats. The findings justified the antidiabetic effect of H. cordata which is attributed to an upregulation of GLUT-4 and potential antioxidant activity, which may play beneficial role in resolving complication associated with diabetes.
Introduction
Diabetes mellitus (DM) is a disease that results in chronic inflammation and apoptosis in pancreatic islets in patients with either type 1 or 2 DM and is characterized by abnormal insulin secretion [1]. Insulin-resistant glucose use in peripheral tissues such as muscle and adipose tissues is a universal feature of both insulin-dependent DM and noninsulin-dependent DM. In this process, glucose transporters (GLUTs) play crucial role [2]. Glucose transporter 4 is mainly expressed in skeletal muscle, heart, and adipose tissues which plays critical role in insulin stimulated glucose transport in these tissues, with glucose uptake occurring when insulin stimulates the translocation of GLUT-4 from the intracellular pool to the plasma membrane [3]. Glucose transporter 2, being the primary GLUT isoform in the liver, plays a pivotal role in glucose homeostasis by mediating bidirectional transport of glucose [4]. It is reported that oxidative stress plays a major role in the development of diabetes associated disorders, possibly due to overproduction of reactive oxygen species (ROS) [5]. Glucose and lipid metabolism are largely dependent on the mitochondrial functional state and physiology which, on excessive ROS formation, leads to mitochondrial oxidative damage and reduced mitochondria biogenesis that contributes to insulin resistance and associated diabetic complications [6,7].
Medicinal plants continue to be an important source in search of a suitable active principle(s), wherein they are currently being investigated for their potential pharmacological properties in the regulation of conditions such as elevated blood glucose level in diabetes [8]. Houttuynia cordata Thunb. (HC) is a single species of its genus and is native to Japan, South-East Asia, and Himalayas. Ethnomedically, whole plant of H. cordata is being used for the treatment of diabetes. In the Ri-Bhoi district of Meghalaya, India, whole plant of H. cordata is eaten raw as a medicinal salad for lowering the blood sugar level and is commonly known by the name Jamyrdoh [9,10]. The plant is also used as an ingredient in insulin secretion promoter compositions [11]. In southern China, green leaves and young roots are 2 Advances in Pharmacological Sciences used as vegetable while dry leaves are used to prepare drink by boiling decoction [12,13]. Reported pharmacological activities of plant including hypoglycaemic [14], antileukemic [15], anticancer [16], adjuvanticity [14], antioxidant [17] and inhibitory effects on anaphylactic reaction and mast cell activation [16].
A recent study has shown that the volatile oil from H. cordata restored the alterations in blood glucose, insulin, adiponectin, and connective tissue growth factor levels in diabetic rats, induced by the combination of a high-carbohydrate and high-fat diet, and STZ injection which may be attributed to the reduced insulin resistance, adiponectin, and connective tissue growth factor levels [18]. Jang et al. [19] reported the potential advanced glycation end products formation and rat lens aldose reductase inhibitory activity of two flavonol rhamnosides 4 and 5 isolated from the whole plant of H. cordata. On the basis of the above reports, the present study was undertaken for the first time to assess the mechanism involved in protective role of H. cordata against STZ-induced glucose toxicity using rat liver, pancreas, and adipose tissue as the working model. In addition, a mechanistic approach of H. cordata against STZ-induced inflammation, apoptosis, and mitochondrial dysfunction was proposed for evaluation. Moreover, GLUT-2 in liver and pancreas and GLUT-4 in adipose tissue were expressed to explain the probable mechanism of H. cordata against STZ-induced impaired glucose utilization.
Animals.
Albino rats of Charles foster strain with body weights of (160-200 g) were obtained from the Central Animal House (Registration number: 542/02/ab/CPCSEA), Institute of Medical Science (IMS), Banaras Hindu University (BHU), Varanasi, India. Before and during the experiment, rats were fed with normal laboratory pellet diet (Hindustan lever Ltd., India) and water ad libitum. After randomization into various groups, the rats were allowed to acclimatize for a period of 2-3 days in the new environment before initiation of experiment. The experimental protocol has been approved by the institutional animal ethics committee (Reference number Dean/10-11/58 dated 07.03.2011).
Phytochemical Analysis.
The extract was subjected to various phytochemicals tests to determine the active constituents present in the crude ethanolic of H. cordata [20]. Total phenolic and tannin content in H. cordata was estimated according to the method of Makkar, [21] using Folin ciocalteu reagent, whereas the method proposed by Kumaran and Joel Karunakaran [22] was followed to estimate total flavonoid and flavonol contents in H. cordata.
Further, H. cordata was standardized with quercetin using high performance thin layer chromatography (HPTLC). A stock solution of both H. cordata and standard quercetin in methanol was prepared in concentration of 5 mg/mL and 0.2 mg/mL respectively. Mobile phase for developing the chromatogram was composed of chloroform: methanol and formic acid mixture in the ratio 7.5 : 1.5 : 1 (v/v/v). The study was carried out using Camag-HPTLC instrumentation equipped with Linomat V sample applicator, Camag TLC scanner 3, Camag TLC visualizer, and WINCATS 4 software for data interpretation. The values were recorded and the developed plate was screened and photo-documented at ultraviolet range with wavelength ( max ) of 254 nm.
Oral Toxicity Studies.
Acute oral toxicity study of ethanolic extract from H. cordata was done according to "Organization for Environmental Control Development" guidelines (OECD: Guidelines 425; Up and Down Procedure). The study was performed on 24 h fasted rats by single dose administration each of 2000 and 5000 mg/kg, (p.o.). The toxicity sign and symptoms or any abnormalities associated with the ethanolic extract of H. cordata were observed at 0, 30, 60, 180, and 240 min and then once a day for the next 14 days. The number of rats that survived was recorded at the end of the study period. [23]. The animals were allowed free access to 5% glucose solution to overcome the drug induced hypoglycemia. Diabetes was confirmed after 48 h and then on the 7th day of streptozotocin injection, the blood samples were collected through retroorbital venous plexus under light anesthesia and plasma glucose levels were estimated by enzymatic GOD-PAP (glucose oxidase peroxidase) diagnostic kit method. The rats having fasting plasma glucose (FPG) levels more than 200 mg/dL were selected and used for the present study [24].
Experimental
Design. The diabetic animals were divided into six groups ( = 6). Group-I, normal control (untreated) rats; group-II, diabetic control rats; group-III, diabetic rats given glibenclamide 10 mg/kg orally for 21 days; group-IV, group-V, and group-VI, diabetic rats that received H. cordata extract at 100, 200, and 400 mg/kg, p.o. body weight, respectively, once daily for 21 days. At the 0th, 7th, 14th, and 21st days blood from each rat was collected through retroorbital venous plexus under light anesthesia. Plasma was separated and the FPG level was estimated. Plasma lipid profile (TC, TG, LDL, HDL, and VLDL), insulin, and other biochemical parameters, that is, creatinine (CRT), blood urea nitrogen (BUN), and total protein (TPR), were also estimated on the 21st day of the experiment.
Evaluation of Mitochondrial Function and Oxidative Stress
2.8.1. Mitochondria Isolation Procedure. Mitochondria were isolated by standard differential centrifugation [25]. The liver, pancreas, and adipose tissue were homogenized in (1 : 10, w/v) ice cold isolation buffer (250 mM sucrose, 1 mM EGTA, and 10 mM HEPES-KOH, pH 7.2). Homogenates were centrifuged at 600 ×g/5 min and the resulting supernatant was centrifuged at 10,000 ×g/15 min and supernatant discarded. Pellets were next suspended in medium (1 mL) consisting of 250 mM sucrose, 0.3 mM EGTA, and 10 mM HEPES-KOH, pH 7.2 and again centrifuged at 14,000 ×g/10 min. All centrifugation procedures were performed at 4 ∘ C. The final mitochondrial pellet was resuspended in medium (1 mL) containing 250 mM sucrose and 10 mM HEPES-KOH, pH 7.2, and used within 3 h. Mitochondrial protein content was estimated using the method of Lowry et al. [26].
Estimation of Mitochondrial Antioxidant Enzymes.
Mitochondrial malondialdehyde (MDA) content was measured based on the TBA reaction test [27]. The activity of superoxide dismutase (SOD) was assayed by the method of Kakkar et al. based on the formation of NADH-phenazine methosulphate-nitro blue tetrazolium formazan measured at 560 nm against butanol as blank [28]. Decomposition of hydrogen peroxide in presence of catalase (CAT) was followed at 240 nm [29]. The results were expressed as units (U) of CAT activity/min/mg of protein.
Estimation of Mitochondrial Succinate Dehydrogenase
Activity (SDH). The mitochondrial succinate: acceptor oxidoreductase (EC 1.3.99.1) was determined by standard protocol based on the progressive reduction of NBT to an insoluble colored compound [a diformazan (dfz)] used as a reaction indicator [30]. The reaction of NBT was mediated by H + released in the conversion of succinate to fumarate. The concentration of NBT-dfz produced was measured at 570 nm. The mean SDH activity of each region was expressed as micromole formazan produced per min per microgram of protein. (MMP). The Rhodamine dye taken up by healthy mitochondria was measured by fluorometric methods [31]. The mitochondrial suspension was mixed with TMRM solution. The mixture was then incubated for 5 min at 25 ∘ C temperature and any unbound TMRM was removed by frequent washings (four times). Then the buffer was added to make up the final volume and florescence emission was read at an excitation 535 ± 10 nm and emission of 580 ± 10 nm using slit number 10. The peak fluorescence intensity recorded was around 570 ± 5 nm. The results are expressed as fluorescence intensity value per milligram of protein.
Pancreatic Histology.
For histopathological studies the pancreas was blotted, dried, and fixed in 10% formalin for 48 h. Thereafter, the tissues were dehydrated in acetone for 1 h and embedded in paraffin wax. Section of pancreatic tissues was then taken through microtome and stained with Haematoxylin-Eosin for photomicroscopic observation [32], which was carried out on Nikon Trinocular Microscope, Model E-200, Japan. analysed by one-way ANOVA, followed by Tukey's multiple comparison test. Data are expressed as mean ± SEM. A level of < 0.05 was accepted as statistically significant.
Phytochemical Analysis.
Preliminary phytochemical analysis of the extract revealed the presence of phenols, flavonoids, tannins, alkaloids, steroids, and carbohydrates as a major component. Total phenolic content of H. cordata was reported to be 45.74 mg/g gallic acid equivalent while total tannin content was estimated to be 33.29 mg/g tannic acid equivalent. Total flavonoid and flavonol contents were found to be 104.55 and 17.16 mg/g rutin equivalents. HPTLC studies revealed well-resolved peaks of H. cordata containing quercetin. The spots of the entire chromatogram were visualized under UV 254 nm and the percentage of quercetin ( 0.51) in H. cordata extract was found to be 4.39% (w/w) ( Figure 1).
Effects of H. cordata on FPG Levels in STZ-Induced
Diabetic Rats. (Table 3). Table 4 represents the effect of H. cordata on body weight of treated rats. Although the mean body weight of treated groups (100 and 200 mg/kg; p.o.) was higher than diabetic control group, it was not statistically significant. However, the rats treated with glibenclamide (10 mg/kg; p.o.) and H. cordata (400 mg/kg, p.o.) showed significant increase in body weight compared to diabetic control. (Table 5). Post hoc analysis revealed that hyperglycaemia significantly increased MDA levels compared to normal control. However, Values are expressed as mean ± SEM of 6 animals in each group. One-way ANOVA showed a significant difference in drug treatment between the groups for HC for total cholesterol, triglyceride, very low density lipoprotein (VLDL), HDL-cholesterol (HDL-C), and low density lipoprotein (LDL). a < 0.05 compared to normal control; b < 0.05 compared to diabetic control; c < 0.05 compared to glibenclamide; d < 0.05 compared to HC 100; e < 0.05 compared to HC 200 (one-way ANOVA followed by Tukey's multiple comparison test). Values are mean ± SEM of 6 animals in each group. One-way ANOVA showed a significant difference in drug treatment between the groups for HC for total creatinine (CRTN), blood urea nitrogen (BUN), and total protein (TPR); a < 0.05 compared to normal control; b < 0.05 compared to diabetic control; c < 0.05 compared to glibenclamide; d < 0.05 compared to HC 100; e < 0.05 compared to HC 200 (one-way ANOVA followed by Tukey's multiple comparison test). Values are mean ± SEM of 6 animals in each group. One-way ANOVA reveals that there were significant differences among the experimental groups [ (5, 30) = 14.12, < 0.05]. a < 0.05 compared to normal control; b < 0.05 compared to diabetic control; c < 0.05 compared to glibenclamide. (One-way ANOVA followed by Tukey's multiple comparison test.) 6 Advances in Pharmacological Sciences
Effect of H. cordata Extract on Mitochondrial Function.
The mitochondrial function in terms of mitochondrial SDH activity was determined in STZ-induced diabetic animals (Table 5)
Effect of H. cordata Extract on Mitochondrial Membrane
Potential (ΔΨ ). The changes in ΔΨ as a marker of mitochondrial integrity during the hyperglycaemic condition are represented in Figure 3. One-way ANOVA showed that there
PPAR-Expressions in Liver, Pancreas, and Adipose
Tissue. Figure 5 shows PPAR-expressions as a marker of inflammation in all three regions of normal control, diabetic control, and diabetic rats subjected to glibenclamide and H. cordata treatment after 21 days. The level of PPARexpression was significantly increased among the groups in the pancreas [ (5, 12) = 6.51, < 0.005]; however, there was no significant change observed in liver [ (5, 12) = 1.93, < 0.005] and adipose tissue [ (5, 12) = 0.79, < 0.005] compared to normal control rats. Further, post hoc analyses revealed that glibenclamide and H. cordata had no significant effect on PPAR-expression.
Effect of H. cordata on GLUT-2 in Liver and Pancreas and GLUT-4 in Adipose Tissue.
As there was an increase in plasma insulin levels in H. cordata-treated diabetic rats and because of the physiologic importance of insulin-dependent GLUT-2 and GLUT-4 translocation to the cell membrane, attempts have been made to see the effect of H. cordata treatment on GLUT-4 level in adipose tissue membrane and GLUT-2 levels in liver and pancreas. In the liver and adipose tissue membrane fractions of diabetic rats, the translocation of GLUT-2 and GLUT-4 was very much reduced when compared with the band density of normal controls. This is quite rational because the deficiency of insulin in the diabetic state would decrease the translocation of GLUT-2 and GLUT-4 from the vesicles to cell membranes. Treatment with H. cordata resulted in the significant increase in membrane GLUT-2 and GLUT-4 levels at the dose of 400 mg/kg. However, there was no significant effect on GLUT-2 level in pancreatic cells. The modulation of GLUT-4 and GLUT-2 protein could thus be one of the mechanisms of antidiabetic properties of H. cordata (Figures 6 and 7).
Histopathological Studies.
The effects of H. cordata on pancreatic cells are represented in Figure 8. The pancreas of the normal rats showed normal islets with intact -cells, whereas, in case of diabetic control rats, atrophy of -cells with vascular degeneration in islets was observed. The rats treated with glibenclamide (10 mg/kg; p.o.) and H. cordata (400 mg/kg; p.o.) depicted regeneration of -cells which were found to be intact and also preserved islets justifying its protective effect.
Discussion
Regular administration of ethanolic extract of H. cordata for 3 weeks resulted in a significant diminution of FPG level with respect to diabetic rat, which clearly explains its antidiabetic activity. The results demonstrated a dose-dependent effect of H. cordata treatment in decreasing FPG. Treatment with H. cordata (200 and 400 mg/kg, p.o.) not only lowered the TC, TG, and LDL level, but also enhanced the HDL-cholesterol which is known to play an important role in the transport of cholesterol from peripheral cells to the liver by a pathway termed "reverse cholesterol transport, " and is considered to be a cardioprotective lipid [33]. Decreased levels of BUN, creatinine and elevation in total protein again indicated that H. cordata can improve renal and liver function [34]. Dysfunctional mitochondria produce excessive amounts of ROS such as superoxide (O 2 − ), hydrogen peroxide (H 2 O 2 ), and peroxynitrite (ONOO − ). This, over production of ROS accumulated in the mitochondrial matrix, leads to collapse of mitochondrial membrane potential (ΔΨ ), decrease in ATP production, and subsequent mitochondrial dysfunction [35]. In line with earlier studies, we also observed that STZ administration produced an increase in the oxidative damage and decreased the antioxidant enzyme activity [36]. Further, there was a significant decrease in mitochondrial function and integrity with the administration of STZ. This effect has been observed in other studies also [37]. H. cordata extract attenuated the STZ-induced mitochondrial oxidative stress and stabilized the mitochondrial function and integrity in the pancreatic tissues. It is well accepted that PPAR-plays a significant role in the pathogenesis of inflammation in several tissues [38]. Moderate amounts of PPAR-are expressed in pancreatic -cells, which increases in the diabetic state [39], leading to accumulation of intracellular triglyceride. In the present study, the level of expression of PPAR-was elevated in STZ-induced rats similar to earlier reports. H. cordata extract did not cause any change in the STZ-induced inflammation indicating that the extract was probably ineffective against STZ-induced inflammation. It is reported that STZ injection causes apoptosis in several tissues such as liver, pancreas, and adipose tissues [37]. As a marker of apoptosis, the level of caspase-3 was increased with the STZ injection in all the tissues under the investigation. H. cordata extract showed significant lowering of caspase-3 level in pancreatic tissues of STZ-injected rats, indicating its promising effect on STZ-induced apoptosis. It is well documented that caspase-3 is a common product of both extrinsic and intrinsic mediated apoptotic pathways [40]. The effect was also well supported by the histopathological studies showing a considerable regeneration in the -cells of pancreas in rats treated with H. cordata 400 mg/kg; p.o. In this context, it is quite impossible to explain the mechanism of H. cordata extract in the STZ-induced apoptosis; however further studies may elaborate the plausible antiapoptotic mechanism of H. cordata extract in STZ-induced model. Several tissues are involved in maintaining glucose homeostasis. Among them, liver, pancreatic -cells, and adipose tissue are the most important because they can sense and respond to changing blood glucose levels. Glucose is taken up into the cell through GLUT-2 and GLUT-4 in the plasma membrane of the cells. In pancreatic -cells, glucose is the primary physiological stimulus for insulin secretion [41]. GLUT-2 is known to play more permissive roles, allowing rapid equilibration of glucose across the plasma membrane. However, it is also essential in glucose stimulating insulin signal (GSIS) because normal glucose uptake and subsequent metabolic signaling for GSIS cannot be achieved without GLUT-2. In diabetic subjects, GLUT-2 and GLUT-4 expression is decreased before the loss of GSIS [42]. Our study suggests that the modulation of GLUT-2 and GLUT-4 protein could thus be one of the mechanisms of antidiabetic potential of H. cordata.
The study revealed a significant increase in serum insulin level in rats treated with H. cordata (especially at 400 mg/kg, p.o.) and glibenclamide as a result of regeneration of pancreatic -cells which were destroyed by streptozotocin .05 compared to normal control; b < 0.05 compared to diabetic control; c < 0.05 compared to glibenclamide; d < 0.05 compared to HC 100; e < 0.05 compared to HC 200 (one-way ANOVA followed by Tukey's multiple comparison test). [43]. Thus, the antidiabetic effect of H. cordata could be attributed to upregulation of GLUT-2 and GLUT-4 protein expressions resulting in potentiation of pancreatic secretion of insulin from existing -cells of islets. Moreover, the study also demonstrated the beneficial role of H. cordata in attenuating oxidative stress responsible for mitochondrial dysfunction. Many works in the literature have shown the antioxidative, anticarcinogenic, antimicrobial, antidiabetic, and antiinflammatory activities of phenols, flavonoids, and polysaccharides [44]. Among the polyphenols, gallic acid, resveratrol, and quercetin are widely distributed in the plant kingdom and are reported to possess antioxidant and antidiabetic properties [45]. In contrast to our results, H. cordata shows anti-inflammatory activity in diabetic condition by improving the level of adiponectins [46]. This discrepancy in the results could be due to the different sets of diabetic condition. The above experiment was performed in the cell lines; however the present study was investigated as an in vivo model of diabetes. Moreover, in both studies H. cordata showed antiapoptotic effect in pancreatic tissues [46]. Studies on diabetic animal models have shown that quercetin significantly decreases the blood glucose level, plasma cholesterol, and TG in diabetic rats, in dose-dependent manner [47]. Beneficial effects of quercetin in increasing the number of pancreatic islets and protective effect on degeneration of -cells along with facilitation in translocation of GLUT-4 have also been reported in the literature [48,49]. Thus, the presence of quercetin quantified in H. cordata may play a contributing factor to the observed antidiabetic activity via the above mentioned pathways.
Conclusion
In conclusion, the present study justified the protective role of H. cordata on pancreatic -cells under high glucose toxic condition by reducing ROS-induced oxidative stress and apoptosis. These findings demonstrated that H. cordata can be employed as a potential pharmaceutical agent against glucotoxicity induced by hyperglycaemia and in oxidative stress associated with diabetes.
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The authors alone are responsible for the content and writing of the paper. | 5,192.6 | 2014-02-24T00:00:00.000 | [
"Biology"
] |
Leveraging Syntactic Constructions for Metaphor Identification
Identification of metaphoric language in text is critical for generating effective semantic representations for natural language understanding. Computational approaches to metaphor identification have largely relied on heuristic based models or feature-based machine learning, using hand-crafted lexical resources coupled with basic syntactic information. However, recent work has shown the predictive power of syntactic constructions in determining metaphoric source and target domains (Sullivan 2013). Our work intends to explore syntactic constructions and their relation to metaphoric language. We undertake a corpus-based analysis of predicate-argument constructions and their metaphoric properties, and attempt to effectively represent syntactic constructions as features for metaphor processing, both in identifying source and target domains and in distinguishing metaphoric words from non-metaphoric.
Metaphor Background
Metaphor can be understood as the conceptualization of one entity using another. Lakoff and Johnson's seminal work shows that metaphors are present at the cognitive level and expressed linguistically (Lakoff and Johnson, 1980). A typical conceptual metaphor mapping is ARGUMENT IS WAR, in which ARGUMENT is structured through the domain of WAR: 1. He defended his position through his publications.
2. Her speech attacked his viewpoint.
The term "linguistic metaphor" is used to indicate these types of words and phrases. We will focus on linguistic metaphor, as identifying these utterances as metaphoric is critical for generating correct semantic interpretations. For instance, in the examples above, literal semantic interpretations of 'defend' and 'attack' will yield nonsensical utterances: a physical position cannot reasonably be defended by a publication, nor can a speech physically attack any kind of entity.
Automatic metaphor processing tends to involve two main tasks: identifying which words are being used metaphorically (here called metaphor identification), and attempting to provide an accurate semantic interpretation for an utterance (here called metaphor interpretation). The first has largely been approached as a supervised machine learning problem, typically using lexical semantic features and their interaction with context to learn the kinds of situations where lexical metaphors appear. The problem of metaphor interpretation is more complex, with approaches including the implementation of full metaphoric interpretation systems (Martin, 1990), (Ovchinnikova et al., 2014), identification of source and target domains (Dodge et al., 2015), developing knowledge bases (Gordon et al., 2015), and providing literal paraphrases to metaphoric phrases (Shutova, 2010), (Shutova, 2013).
In both identification and interpretation systems, syntax tends to play a limited role. Many systems rely only on lexical semantics of target words, or use only minimal context or dependency relations to help disambiguate in context (Gargett and Barnden, 2015), (Rai et al., 2016). Others rely on topic modeling and other document and sentence level features to provide general semantics, and compare the lexical semantics to that, ignoring the more "middle"-level syntactic interactions (Heintz et al., 2013). While these approaches have been effective in many areas, there is evidence that figurative language is significantly influenced by syntactic constructions, and thus if they can be represented more effectively, metaphor processing capabilities can be improved.
We will examine five kinds of predicateargument constructions in corpus data to assess their metaphoric distributions and usefulness as features for classification. Our contribution is twofold. First, we examine the LCC metaphor corpus, which includes source and target annotations, to determine their use in predicateargument constructions (Mohler et al., 2016), and employ syntactic representations as features to improve source/target classification. Second, we investigate predicate-argument constructions in the VUAMC corpus of metaphor annotation (Pragglejaz Group, 2007), and employ syntactic features to predict metaphoric vs non-metaphoric words.
Metaphor and Constructions
Recent metaphor research has indicated that construction grammar can be employed to determine the source and target domains of linguistic metaphors (Sullivan, 2013). In many cases, certain constructions can determine what syntactic components are allowable as source and target domains. For example, verbs tend to evoke source domains. The target domain is then evoked by one or more of the verb's arguments (from Sullivan pg 88): 1. the cinema beckoned (intransitive) 2. the criticism stung him (transitive) 3. Meredith flung him an eager glance (ditransitive) In these instances, the verb is from the source domain and at least one of the objects is from the target. However, arguments can also be neutral and don't necessarily evoke the target domain. Pronouns like 'him' in (2) and (3) don't evoke any domain. The optionality of domain evocation makes it harder to predict which elements of the construction participate in the metaphor. Despite this limitation, this analysis shows that syntactic structures beyond the lexical level can be indicative of source and target domains. To better understand how these structures determine metaphor, we explored metaphor-annotated corpus data for predicate-argument constructions.
Computational Approaches
While metaphor processing has largely been focused on capturing lexical semantics, there have been a variety of approaches that incorporate syntactic information. Many computational approaches focus on specific constructions, perhaps indicating the need to classify different metaphoric constructions through different means. The dataset of (Tsvetkov et al., 2014) provides adjective-noun annotation which has been extensively studied (Rei et al., 2017), . A particularly promising approach is that of , who use compositional distributional semantic models (CDSMs) to represent metaphors as transformations in vector space, specifically for adjective-noun constructions. Another relevant approach is that of (Haagsma and Bjerva, 2016) who use clustering and selectional preference information to detect metaphors in predicate argument constructions, including verbs with objects, subjects, and both. Their highest F1 is 57.8 for verbs with both arguments.
Many systems that rely heavily on lexical resources also include some dependency information. (Rai et al., 2016) and (Gargett and Barnden, 2015) use a variety of syntactic features including lemma, part of speech, and dependency relations. However, both systems are feature-rich and these syntactic elements' contribution is unclear. (?) use lexical features along with contrasting those features between the target word and its head. (Dodge et al., 2015) employ a variety of constructions in identifying metaphoric source and target domains. They identify a broad range of constructions and use these as templates that metaphoric expressions can fill. Our work expands on this idea by formalizing the constructions into features for statistical metaphor identification.
Perhaps the most syntactically oriented metaphor identification system is that of (Hovy et al., 2013), who uses syntactic tree kernels to identify metaphor. They use combinations of syntactic features via tree kernels and semantics via WordNet supersenses and target word embeddings. Our approach expands on this by exploring different syntactic representations and incorporating semantics through word embeddings into the syntactic structures.
Corpus Analysis
Sullivan identifies a large number of constructions and the possible configurations of their arguments with regard to source and target domains. While some corpus examples are provided that show the variety of source-target patterns in each construction's argument structure, an in-depth analysis of how these constructions and their metaphoric properties are distributed is still needed. We examined the predicate argument constructions they analyze by using hand-annotated metaphor corpora to better understand the distributional patterns that occur. This allows us to make predictions about what kind of constructions and arguments are useful for metaphor identification and interpretation and what might be a computationally feasible way to implement them.
While they examine many kinds of constructions, most of them seem based almost entirely on the lexical semantics of the words involved, and thus can be captured simply by effectively representing the meaning of individual words. Domain and predicative adjective constructions fall into this category: the construction is identified by the type of adjective, which needs to be represented at the lexical level. The more interesting cases are argument structure constructions, which take many forms. Sullivan identifies nine different argument structure constructions that each have their own source and target properties: To identify the use of metaphor in these constructions, we will rely on two resources: the LCC metaphor corpus and the VUAMC corpus. The freely available portion of the LCC corpus contains approximately 7,500 source/target pairs, allowing for a more in-depth look at metaphoric semantics. The VUAMC contains approximately 200,000 words of text with each word tagged as metaphoric or non-metaphoric. This allows for large scale analysis of metaphoricity versus nonmetaphoricity at the word level.
Identifying Constructions
To examine metaphors in these corpora, we need a method for automatically identifying predicateargument constructions. The VUAMC corpus, as a subsection of the BNC baby, comes with gold-standard dependency parses. For the LCC dataset, we used the dependency parser from Stanford Core NLP tools . These parses are sufficient to identify intransitives, transitives, and ditransitive constructions. Verb instances that have an indirect object are ditransitive, those that lack an indirect object but have a direct object are transitive, and those that lack either but have a subject are intransitive. Copulas are marked in the dependency parses, so we can easily identify equative constructions. While similes can take many forms, Sullivan's work focuses on simile constructions that consist of a copular verb and the word 'like'. This oversimplifies to some degree, as many similes don't need a copula ('she fretted like a mother hen', 'they flew like bats'), but it allows us to create a subset of equative constructions that represent copular similes.
This analysis is necessarily limited, as the we cannot automatically capture more complex constructions via dependency parses, and many of these are often metaphorically rich. While we understand this limitation, we believe that we can utilize syntactic features of these basic constructions as a starting point, with a future goal of expanding to more complex examples.
Also note that we only identify the surface realization of these constructions -any dropped arguments or missing elements that aren't in the dependency parse aren't considered a part of the construction. Thus we see examples of typically ditransitive verbs (like 'give') that occur intransitively and transitively, as they lack overt direct and indirect objects.
LCC Analysis
To explore source and target domains, we employ the free portion of the LCC corpus from Mohler et al, which contains approximately 7,500 source/target metaphor pairs in sentential context, rated from 0 to 3 on their degree of metaphoricity. For our research, we included only those instances that were rated above 1.5, yielding approximately 3,000 metaphoric sentences. These annotations also include the source and target domains of the metaphors, and the lexical trigger phrases that engender the source and target domains. This allows us to quantify Sullivan's analysis of source and target domains in different constructions, and shows the actual distribution of source and target domain items in each construction.
In order to identify constructions in the LCC data, we extracted syntactic relations from the dependency parses, using the basic patterns previously defined to identify predicate argument constructions. This allows us to identify the five different constructions: intransitives, transitives, ditransitives, equatives (copulas), and similes (analyzed as a subset of equative constructions). For each construction found, we can identify the predicate and the predicate's arguments, and determine for each whether they are identified as metaphoric and whether they belong to the source or target domain. The vast majority of constructions in the LCC are intransitive, transitive, and equative. Ditransitives (.4%) and similes (.1%) are exceedingly rare. This may be because the similes found are only the verbal type: instances of a copula with the word 'like'. Other similes are likely missed by this automatic approach.
The majority of metaphoric verbs (92%) are source domain items, supporting Sullivan's claims. Subjects and objects tend to be from the target domain (61% each). Ditransitive verb constructions are relatively rare, with only 43 found, and only 3 of those containing a metaphoric verb. Figure 1 shows the counts of source and target items in the LCC data, based on construction and argument of the construction. Note that in equative constructions, direct objects are almost always source domain items, showing a parallel between copular arguments and verbs. This is likely due to the predicative nature of the direct objects of copular verbs.
Source and Target Identification
Given that verbs and their argument structures have varying distributions of source and target domain items, we believe that these syntactic structures can be effectively employed in the classification of source and target domain words. While identifying source and target domains at the sentence level requires lexical and sentential semantics and may not require syntactic information, identifying lexical triggers can be improved by using better syntactic representations. To this end we set up a classification task for identifying source and target elements. The LCC contains phrase-level annotations for source and target elements. We split each sentence into words, projecting the source and target annotations to the word level. From this, we developed three classification tasks: (1) identifying source words, (2) identifying target words, and (3) identifying any metaphoric word (either source or target). Our classification scheme focuses on verbs and nouns, as these are the elements that compose the syntactic structures in question.
We developed a set of different representations designed to capture construction-like structures, and employ them for source/target classification. This approach follows the intuition of (Hovy et al., 2013): "metaphorical use differs from literal use in certain syntactic patterns". We implemented this theory by developing various representations of constructional syntax and pairing them with lexical semantic features.
For our lexical semantics component, we experimented with the word embeddings from word2vec (Mikolov et al., 2013), using the pretrained Google News data, as well as the Glove embeddings (Pennington et al., 2014). We found in validation that the Google News vectors yielded slightly better performance, and so those were used in further experiments.
Predicate Argument Construction
For a basic integration of syntax, we used the above corpus analysis technique to identify which predicate-argument construction the verb token belongs to. This results in a one-hot vector representing either an intransitive, transitive, ditransitive, equative, or simile construction. This provides basic, purely syntactic knowledge of how many arguments this particular instance of a verb currently has. For nouns, we extend this to include which slot in the construction the noun is filling (subject, direct object, indirect object) in addition to the type of predicate-argument construction.
Head and Dependent Features
Including representations of the head word and dependent words of the word to be classified is a straightforward way to include basic syntactic information. For verbs, this mainly involves the dependents, although many verbs also have head words. We include a concatenation of the average embedding over the word's dependents and the embedding of the word's head.
Dependency Relations
A more general and perhaps more powerful way of converting dependency relations into syntactically relevant features is to include the specific dependency relations for each dependent of the target. For verbs, these include things like subjects, direct objects, adverbial modifiers, nominal modifiers, passive subjects, and more. Capturing the fine-grained dependencies for each verb is analogous to determining the exact syntactic construction it is being realized in. Combining this feature with the embeddings of dependents and heads is a promising avenue for linking syntax and semantics.
VerbNet Class
VerbNet is a lexical semantic resource that groups verbs into classes based on their syntactic behavior (Kipper-Schuler, 2005). It categorizes over 6,000 verbs into classes, each of which contains syntactic frames that the verbs in the class can appear in. It also contains distinct senses, allowing it to distinguish between different verb uses in context. Previous approaches have employed VerbNet as a lexical resource (Beigman Klebanov et al., 2016), but aggregated the senses of each verb, removing the syntactic distinctions that VerbNet makes for different word senses. We ran word-sense disambiguation to determine the VerbNet class for each verb token (Palmer et al., 2017). We included one-hot vectors representing verb senses for each token, and combining this with knowledge of the particular constructions and the lexical semantics provided by embeddings for each token gives syntactically motivated information about the semantics of the utterance. For noun identification, we include the VerbNet class of the head of that noun.
Experiments
As a baseline, we began with using the embedding of the word to be classified. We concatenated this with the embeddings of the single previous and following words, as this proved the best context in our validation. This creates a representation of lexical semantics and a word's context, without any specific knowledge of the syntactic relations the word is involved in. We then added each syntactic representation. These experiments were done using a training-validation-test split of 76/12/12. We experimented with Maximum Entropy, Naive Bayes, Random Forest and Support Vector Machine classifiers, and through validation chose a SVM with a linear kernel, L2 regularization and squared hinge loss. We then ran the classifier using our baseline, and added each feature separately. Finally, we combined the best feature set for each classification task, judged by the improved performance of each feature over the baseline. The classification was split into three tasks: identifying source items, identifying target items, and identifying metaphoric (either source or target) from non-metaphoric. The results of these experiments are in table 2.
From these results we can see that classifying source-domain words in the LCC data is harder than classifying target-domain words. This may be because of the broad range of domains, as the corpus contains 114 possible source domains. Target items are much easier to classify, likely because the dataset contains only a limited number (32) of target domains. Embeddings are effective at representing semantics, and they can accurately determine the domain of lexical items, allowing for easy classification of target items.
Our syntactic features show mixed results. Adding sentential context is consistently effective, showing that naive contextual approaches are helpful. Adding dependency embeddings is also consistently effective, supporting our hypothesis that knowledge of syntactic properties can be helpful in metaphor classification. Other syntactic features are inconsistent, especially in predicting the metaphoricity of verbs. Selecting only the feature sets that showed improvement over the baseline yields the best results for most categories.
VUAMC Analysis
The LCC allows for an in-depth examination of source and target domains, but is relatively small compared to the VUAMC. We can use the VUAMC data to inspect the distribution of word metaphoricity with regard to argument structure constructions. While Sullivan's work focuses on source and target domain elements and not whether or not words are used metaphorically, we can examine the binary classifications in the VUAMC to provide insight into the distribution of metaphoric verbs and the predicate-argument constructions they participate in. Counts of argument structure verbs and arguments and their metaphoricity are shown in table 3.
From the data in table 3, we can see clear distinctions between different constructions and the metaphoricity of their arguments. Verbs in instransitive constructions are much less likely to be metaphoric than those used in transitives, and both less so than those in ditransitive constructions. The VUAMC chooses not to mark copular verbs as metaphoric, and only one instance was found of equative constructions having a metaphoric verb.
We might expect that different constructions would also impact the distribution of the predicates' arguments. However, from the data we see that verb arguments are fairly consistent. Indirect objects in ditransitive constructions were never observed to be metaphoric, but direct objects are between 11% and 16% metaphoric throughout. Subjects vary from 2.8% in ditransitives to 11.7% in equative constructions. One distinctive feature is that subjects are much less likely than objects to be metaphoric.
The overall distribution of metaphoric uses by verb construction shows that the more arguments that are present in the construction, the more likely the verb is being used metaphorically. For further evidence, we can examine the distribution of metaphoric usages on a verb-specific basis.
We calculated the average metaphoricity of each verb found in the VUAMC, and sorted them by the type of construction they are found in. We performed this analysis on a type and token basis, shown in figures 2 and 3. From the data, we see that the majority of verbs in all constructions are used exclusively non metaphorically. While a large number of verb types only occur metaphorically, this accounts for a much smaller number of verb tokens. Verb types that occur only metaphorically are relatively rare. We can also see that ditransitive and copula verb types are exceedingly rare, but copula tokens are very common and almost always literal. We extended this analysis by examining the distribution of the verb types that can appear intransitively, transitively, and ditransitively. Our hypothesis in studying these verbs is that the type of construction the verb appears in is predictive of that verb's metaphoric use, independent of the verb's overall behavior. Eleven verbs appeared in all three constructions, and the analysis of their metaphoricity is presented in figure 4.
From the distribution in the VUAMC corpus, the data indicates that the type of argument structure construction does not significantly change the distribution of metaphoricity. The verbs generally have the same percentage of metaphoric usages regardless of which construction they appear in. Only 'give' appears in more than 2 instances of the ditransitive, and its distribution mirrors that of its use in other constructions.
Two components from our corpus analysis stand relevant for automatic metaphor processing. First, in broad scope over all verb tokens, predicates' metaphor distributions are dependent on the kind of construction they occur in. Second, the verb itself is critical, as each verb tends to follow the same pattern of metaphoricity throughout its constructions. This supports our belief that identification of metaphor requires modeling of the interaction of syntactic and semantic information.
Metaphor Identification (VUAMC)
We employ the same experimental set up of the previous classification task using the VUAMC. The VUAMC doesn't contain source or target annotations, so the classification problem is limited to identifying metaphoric words from nonmetaphoric words. We employ the same baseline and syntactic representation features. Again, we For metaphoric identification in the VUAMC, all of the syntactic features improved classification over the baseline for verbs. For nouns, the dependency embeddings and VerbNet class of the noun's head were effective. For both, combining all of the syntactic representations yields the best performance. While this classification based on syntactic is slightly lower than some recent experiments (Beigman Klebanov et al., 2016), it shows improvement over using purely lexical semantics, and we believe the incorporation of better syntactic representations can be used to improve metaphor identification systems.
Conclusions
The type of syntactic construction a verb is present in provides unique evidence of how it is being used metaphorically. It is important to effectively inte- grate syntax and semantics to detect and interpret metaphor, and because there are so many types of metaphors and they occur in such a wide array of contexts, it may be helpful to use separate methods of representing metaphoric semantics depending on the syntactic constructions involved. While our results indicate that these integrations of syntactic representations do not yet achieve state of the art performance, we believe that improving representations of syntactic constructions can provide some benefit to metaphor processing. To that end, our future goals include exploring better representations of the interaction between syntax and semantics. Models like syntactic tree kernels, compositional distributional semantic models, and other syntactically driven methods are likely to improve classification if they can properly combine syntactic and semantic representations. Additionally, as different constructions are likely to yield different types of metaphoricity, we aim to employ ensemble methods that incorporate construction-based knowledge to select the most effective classifier, and extending our approach to identifying source and target domains in addition to lexical triggers. | 5,434.8 | 2018-06-01T00:00:00.000 | [
"Linguistics",
"Computer Science"
] |
Is there a breakdown of effective field theory at the horizon of an extremal black hole?
Linear perturbations of extremal black holes exhibit the Aretakis instability, in which higher derivatives of a scalar field grow polynomially with time along the event horizon. This suggests that higher derivative corrections to the classical equations of motion may become large, indicating a breakdown of effective field theory at late time on the event horizon. We investigate whether or not this happens. For extremal Reissner-Nordstrom we argue that, for a large class of theories, general covariance ensures that the higher derivative corrections to the equations of motion appear only in combinations that remain small compared to two derivative terms so effective field theory remains valid. For extremal Kerr, the situation is more complicated since backreaction of the scalar field is not understood even in the two derivative theory. Nevertheless we argue that the effects of the higher derivative terms will be small compared to the two derivative terms as long as the spacetime remains close to extremal Kerr.
Introduction
Extremal black holes (BHs) are an important special class of BHs with degenerate, zero temperature horizons. They play a prominent role in String Theory as they are often supersymmetric and do not evaporate. As distinguished members of the BH family with broad theoretical applications, understanding their classical stability properties seems important. Are extremal BHs classically stable?
While proving the nonlinear stability of the Kerr BH remains as a major goal of mathematical relativity, some significant steps towards this goal have already been made. The current state-of-the-art are the recent proofs of linear stability of Schwarzschild under gravitational perturbations [1] and linear stability of a massless scalar on Kerr [2]. Importantly, these proofs are restricted to non-extremal BHs. The reason is that the so-called horizon redshift effect is essential in those analyses. This is the phenomenon that outgoing radiation propagating along the future event horizon suffers a redshift and therefore decays. The characteristic decay time is proportional to the BH's surface gravity. At extremality the surface gravity vanishes so there is no horizon redshift effect and the stability proofs fail.
The search for a new approach to study the stability of extremal BHs led Aretakis, in a series of works [3,4,5,6], to prove that massless scalar perturbations of extreme Reissner-Nordström (RN) and axisymmetric massless scalar perturbations of extreme Kerr BHs display both stable and unstable properties. He showed that the scalar field and its derivatives decay outside the event horizon. However, on the event horizon, the absence of a horizon redshift effect means that outgoing radiation propagating along the event horizon does not decay. Mathematically, this means that a transverse derivative of the scalar field does not decay along the horizon and higher transverse derivatives grow with time. For spherically symmetric massless scalar perturbations of extreme RN, derivatives blow up at least as fast as Killing time coordinate). In [11] an extension to charged perturbations of extreme RN was discussed; these were shown to resemble non-axisymmetric modes in extreme Kerr.
The above discussion concerns linear perturbations of extreme BHs. It is natural to ask what happens when one considers nonlinearity and backreaction. Aretakis considered the case of a scalar field with a particular kind of self-interaction and found that the nonlinearity made the instability worse, leading to a blow up in finite time along the event horizon [12]. A different kind of nonlinearity was considered in Ref. [13], for which it was found that the nonlinearity did not lead to any qualitative difference from the linear equation. However, for both of these examples, the nonlinearity was not of a kind that would arise in physical applications. The backreaction problem was investigated numerically in Ref. [14]. It was found that, for a generic (massless scalar field) perturbation, an extreme RN black hole will eventually settle down to a non-extreme RN solution. However, during the evolution, there is a long period when derivatives exhibit the behaviour (1.1), confirming that the instability persists when backreaction is included. Furthermore, by fine-tuning the perturbation it can be arranged that the late-time metric approaches extreme RN, in which case the nonlinear solution exhibits the behaviour (1.1) indefinitely.
We now turn to the physical relevance of the Aretakis instability. If fields decay outside the event horizon then why does it matter that higher transverse derivatives blow up on the horizon? One reason is that we expect the classical equations of motion to be corrected by higher derivative terms, as is the case in string theory. If higher derivatives become large on the horizon then it seems likely that the higher derivative terms in the equation of motion will become large [14]. In other words, the Aretakis behaviour suggests a possible breakdown of effective field theory at late time on the event horizon of an (arbitrarily large) extreme black hole. 1 The aim of this paper is to investigate whether or not higher derivative corrections to the equations of motion become important during the Aretakis instability or the even worse non-axisymmetric extremal Kerr instability of Ref. [9]. We will consider a nonlinear theory consisting of Einstein-Maxwell theory coupled to a massless scalar, and then add higher derivative corrections which are restricted only by the requirement of general covariance and a shift symmetry for the scalar field. In section 2 we consider the extremal RN solution. We start with a brief review of the Aretakis instability. We then consider the AdS 2 × S 2 near horizon geometry of an extremal RN black hole, taking into account the higher derivative corrections to the background geometry. We expand on a previous discussion [8] of how the Aretakis instability can be seen in the near-horizon geometry. We then show that, for a large black hole, linear higher derivative corrections lead only to small corrections to Aretakis' results. In particular, the leading (spherically symmetric) instability of the near-horizon geometry is unaffected by these corrections. Ultimately the reason for this is that the higher derivative terms must exhibit general covariance, which implies that they take a very simple form when linearized around a highly symmetric background such as AdS 2 × S 2 .
It is not obvious that this will remain true when we consider the much less symmetric geometry of the full black hole solution. So next we consider the size of (possibly nonlinear) higher derivative terms in all of the equations of motion during the Aretakis instability in the full extreme RN geometry. We argue that such terms remain small compared to the nonlinear 2-derivative terms. Hence there is no indication of any breakdown of effective field theory for extreme RN. Ultimately this result can again be traced back to general covariance restricting the possible form of the higher derivative terms.
In section 3 we discuss the case of extremal Kerr. Again we start by investigating the scalar field instability in the near-horizon geometry. In particular, we give a simple derivation of results analogous to those of Ref. [9] for the scalar field instability in the near-horizon extreme Kerr (NHEK) geometry. We explain how these results are robust against higher derivative corrections of the NHEK geometry. Furthermore, our method can incorporate outgoing radiation at the event horizon in the initial data, unlike the approach of Ref. [9].
Nevertheless, our results are in agreement with those of Ref. [9], indicating that this initial outgoing radiation does not make the dominant (non-axisymmetric) instability any worse. We then consider linear higher derivative corrections to the equation of motion for the scalar field and argue that these just give small corrections to the results, again without making the instability any worse. So, at the level of the near-horizon geometry, there is no sign of any breakdown of effective field theory.
Finally we consider the scalar field instability in the full extreme Kerr geometry. Here the effect of nonlinearities is not yet understood, even in the 2-derivative theory. So we simply assume, in analogy with the nonlinear extreme RN results, that the geometry remains close to extreme Kerr even when 2-derivative nonlinearities are included. With this assumption we estimate the size of higher derivative corrections to the equations of motion. We find that these remain small compared to the 2-derivative terms. So again there is no obvious sign of any breakdown of effective field theory. Once again the reason can be traced to general covariance restricting the form of possible higher derivative terms.
Einstein-Maxwell-scalar theory
Consider an Einstein-Maxwell-scalar theory where the scalar field is massless and minimally coupled. This theory is described by the action 2 .
where F = dA with A a 1-form potential. We now consider higher derivative corrections to this two derivative action. We write the action as where S 2 is as above and where α has dimensions of length and L k is a scalar function of the metric, Maxwell field strength and scalar field, involving k derivatives of the scalar field, metric or electromagnetic potential. We will assume that the scalar field is coupled only through its derivatives so the theory possesses a shift symmetry Φ → Φ + const. Furthermore, we assume that L k does not involve any terms which are linear in (derivatives of) Φ, which implies that setting Φ = const is a consistent truncation of the theory.
Since it is not possible to construct a scalar Lagrangian with 3 derivatives, we have S 3 = 0 and the first higher derivative term in the action is S 4 .
Aretakis instability in 2-derivative theory
First we review the Aretakis instability in the 2-derivative theory. Setting Φ = constant, the two-derivative theory admits the extreme RN black hole as a solution. We write the metric as and the Maxwell field is where dΩ is the volume element on a unit radius S 2 . We have assumed that the black hole is magnetically charged with charge Q. 3 In this background, Aretakis considered linear perturbations in the scalar field, which we write as ψ ≡ δΦ. The equation of motion for ψ in the 2-derivative theory is We can decompose ψ in spherical harmonics: Because of the spherical symmetry we can ignore the dependence on m and just write ψ ℓ . The wave equation becomes Consider first ℓ = 0. Evaluating (2.8) at the horizon δ = 0 shows that the quantity is conserved along the horizon (independent of v), and in particular does not decay, for generic initial data, at late times. H 0 is called an Aretakis constant. Since ψ 0 | horizon itself does decay at late times on the horizon [3], this shows that the first derivative ∂ r ψ 0 | horizon does not decay -instead, it tends to H 0 . Higher derivatives of ψ 0 behave even 'worse' on the horizon: at late times they grow indefinitely, as can be seen by acting on equation (2.8) with ∂ r and restricting to the horizon giving Integrating with respect to v then gives as v → ∞. It follows that This can be extended by induction to an arbitrary number of radial derivatives. Acting with ∂ k−1 r on (2.8), restricting to the horizon and integrating along it, shows that as v → ∞, where here and below we ignore dimensionless constants on the RHS. Hence higher derivatives of ψ 0 grow polynomially with v at late time on the event horizon. This is the Aretakis instability.
Similar behaviour occurs for ℓ > 0. Acting on (2.8) with ∂ ℓ r and restricting to the horizon shows that there is a conserved quantity As in the ℓ = 0 case, an inductive procedure yields, for k ≥ ℓ + 1 at late time along the event horizon. Notice that ℓ + 2 derivatives are required to construct a quantity that grows along the horizon, hence the Aretakis instability is strongest for the ℓ = 0 mode. We will also need to know the behaviour of quantities which decay along the horizon.
Numerical results in Ref. [8] strongly suggest that ψ 0 ∼ v −1−ℓ at least for ℓ = 0, 1. This is confirmed by rigorous results of Ref. [16], which prove that (2.15) holds for any k ≥ 0 when the Aretakis constant H ℓ is non-zero. It is also proved that v-derivatives behave in the way one would expect by naively differentiating w.r.t. v: We have dropped all coefficients on the RHS of (2.16). These coefficients are all proportional to H ℓ multiplied by appropriate powers of Q.
Although the following will not be used in our analysis, it is interesting to note that the above late-time behaviour is reproduced by an expression of the form where f (ℓ) is a smooth function with f (ℓ) (0) = 0. This Ansatz can be substituted into (2.8).
Taking the late time v → ∞ limit, keeping z ≡ vδ fixed, (2.8) then reduces to an ordinary differential equation for f . Solving it gives the 0th order wavefunction (Q = 1): where c i are constants. For ℓ = 0, it reduces to The late time behaviour here involves two constants H 0 and c 20 . The interpretation of the latter is as a Newman-Penrose constant [17]. Just as the Aretakis constants are associated to outgoing radiation propagating along the future event horizon, the NP constants are associated to ingoing radiation propagating along future null infinity. In other words, they correspond to late time ingoing radiation. In equation (2.16) we assumed vanishing NP constants but this result can be generalized to allow non-zero NP constants [16]. Henceforth we will assume vanishing NP constants.
Higher derivative corrections in near horizon geometry
Setting Φ = constant, the two-derivative theory admits the extreme RN black hole as a solution. We assume that this solution can be corrected so that it remains a solution of the theory to all orders in α. We will assume that the corrected black hole is magnetically charged with charge Q defined by (2.5). Of course this satisfies dF = 0.
The near horizon geometry of this black hole will be AdS 2 × S 2 where the AdS 2 and S 2 have radii L 1 and L 2 respectively. We can write L i = QL i (α/Q) i = 1, 2 whereL i is dimensionless. For small α/Q the higher derivative corrections will be negligible and the AdS 2 and S 2 will both have radius Q. The higher derivative corrections start at O(α 2 ) hence we haveL We write the AdS 2 × S 2 metric in ingoing Eddington-Finkelstein coordinates as Ref. [8] showed that a massless scalar in this geometry exhibits the Aretakis instability at the future Poincaré horizon r = 0. At first this seems rather surprising given that a scalar field in AdS 2 × S 2 exhibits no instability in global coordinates. This was discussed in Ref.
[8], we will expand a little on this discussion here. For a well-posed problem we need to impose boundary conditions at infinity in AdS 2 .
Following Ref. [8], we assume that boundary conditions have been chosen such that, in a neighbourhood of r = 0, v → ∞ (where the Poincaré horizon intersects infinity), these conditions correspond to "normalizable" boundary conditions for the scalar field. The Aretakis instability does not involve the growth of some scalar quantity, but is instead associated to the growth of the components of a tensor, specifically the second derivative of ψ. But how does one know that this growth is associated to some physical effect rather than to bad behaviour of the basis in which the components are calculated? The point is that the asymptotically flat black hole solution has a canonically defined Killing vector field V which generates time translations. One can choose a basis to be time-independent, i.e., Lie transported w.r.t. V . If a component of some tensor exhibits growth in such a basis then one can be sure that this is a physical effect rather than an artifact of the choice of basis.
An example of such a basis is a coordinate basis where V is one of the basis vectors. This is the case in Eddington-Finkelstein coordinates where V = ∂/∂v. This is why one can be sure that the Aretakis instability is not a coordinate effect. Since the Aretakis instability can be seen in the near-horizon geometry, we will start by investigating the effect of higher derivative corrections on this instability in the AdS 2 × S 2 background (2.22). We will take into account two sources of higher-derivative corrections: first we are using the exact, higher-derivative corrected, background (2.22). Second, we will include the effect of linear higher derivative corrections to the scalar field equation of motion. The reason for restricting to linear higher derivative corrections is that if we allow nonlinearity then we have to incorporate the effects of the backreaction of the scalar field on the geometry. However, even in the 2-derivative theory, it is known that this backreaction destroys the AdS 2 asymptotics [18]. To incorporate this backreaction we have to consider the full black hole solution, as we will do in the next section.
Since the action does not contain terms linear in Φ, the higher derivative corrections to the Einstein equation and the Maxwell equation also do not contain terms linear in Φ, and the corrections to the scalar equation of motion do not contain any Φ-independent terms. Furthermore, our assumption of a shift symmetry implies that the equations involve only derivatives of Φ. This structure implies that when we linearize around an exact background solution with Φ = const, the linear perturbation to Φ decouples from the linear metric and Maxwell field perturbations.
To discuss linear higher-derivative corrections to the scalar field equation of motion we will work at the level of the action. We expand the action to quadratic order in ψ = δΦ.
We then substitute in the expansion in spherical harmonics (2.7), and perform the integral over S 2 . Modes corresponding to different harmonics will decouple from each other, giving an effective action for the field ψ ℓm in AdS 2 of the form 4 where g 2 is the AdS 2 metric (with radius L 1 ), is the d'Alembertian of this metric, and c ℓn are (real) constants depending on α and Q. The form of this effective action is dictated by the AdS 2 symmetry of the background. Recall our assumption that the scalar field is derivatively coupled. Derivatives can act on either the S 2 or AdS 2 directions. But the spherically symmetric ℓ = 0 mode is constant on S 2 hence it cannot appear without AdS 2 derivatives in the above action. It follows that c 00 = 0. Terms in the action with n ≥ 2 must arise from higher derivative terms in the original action and hence must appear with appropriate powers of α. We can write wherec ℓn is a dimensionless function of α/Q. For n = 0, 1 we can separate out the terms present in the 2-derivative theory from those arising from the higher derivative corrections (to both the background and the equation of motion): 5 Againc ℓn is a dimensionless functions of α/Q andc 00 = 0. A standard result in effective field theory is that the lowest order (i.e. two derivative) equation of motion can be used to simplify the higher derivative terms in the action. This is achieved via a field redefinition [19]. To see how this works here, perform a field redefinition (here we suppress the ℓ, m indices throughout) where the dimensionless coefficients d n (α/Q) are to be determined. We substitute this into the action and let E n be the coefficient ofψ nψ . We demand that E n = 0 for n ≥ 2. This gives a set of equations that can be solved order by order in α/Q to determine the coefficients d n . To lowest order, Plugging the latter back into E 2 = 0 then determines the O(α 2 /Q 2 ) part of d 2 .
One then uses E 4 = 0 to determine d 4 to O(1), plug this back into E 3 = 0 to determine Repeating this process to all orders gives Hence, to all orders in α,ψ ℓm behaves as a massive scalar field in AdS 2 with mass m ℓ . Since ψ ℓm is linearly related toψ ℓm , the same will be true for ψ ℓm . We see that the only effect of the higher derivative corrections is to correct the mass of this scalar field. Of course, all we have done here is to perform a Kaluza-Klein reduction of the scalar field ψ on S 2 .
Note that the higher derivative corrections do not generate a mass for ψ 00 . The masslessness of ψ 00 is protected by the assumed shift symmetry, which implies c 00 = 0 and hence m 2 0 = 0 to all orders. So higher derivative corrections do not change the equation of motion for the ℓ = 0 mode. Now we can discuss the effect of the higher derivative corrections on the Aretekis instability in AdS 2 × S 2 . In the absence of such corrections, this instability is strongest in the ℓ = 0 sector, with ∂ 2 r ψ 00 growing linearly with v along the horizon at r = 0. For higher partial waves more derivatives are required to see the instability: ∂ ℓ+2 r ψ ℓm grows linearly with v. From the results just obtained, we see that higher derivative corrections have no effect on the ℓ = 0 sector and so ∂ 2 r ψ 00 will still grow linearly with v. However, these corrections do affect higher ℓ modes through the change in the mass just discussed. To understand the effect of this change in the mass, we can use results of Ref. [8], which determined the behaviour of massive scalar fields in AdS 2 along the Poincaré horizon at late time. 6 The result is that, for a scalar of mass m, at late time along the horizon r = 0 where ∆ is the conformal dimension with L 1 the AdS 2 radius. So for a massive scalar, ∂ k r ψ decays along the horizon if k < ∆ and grows if k > ∆. Applying this in our case, writing If δM ℓ > 0 then the higher derivative corrections have led to increased stability in the sense that the decay is slightly faster for k < ℓ + 1 and the blow up is slightly slower for k > ℓ + 1.
On the other hand, if δM ℓ < 0 then the higher derivatives lead to reduced stability in the sense that not only do we have faster growth for k > ℓ + 1, we also have power law growth for k = ℓ + 1. In particular, if δM 1 < 0 then the second derivative of the ℓ = 1 mode exhibits power law growth along the horizon. However, the exponent in this power law will be proportional to −δM 1 and therefore small compared to the linear growth exhibited by the second derivative of the ℓ = 0 mode. So even though higher derivative corrections may strengthen the instability in the higher ℓ modes, for small α/Q, they do not strengthen them enough that they compete with the dominant ℓ = 0 mode, which is unaffected by these corrections.
Of course, the question of whether δM ℓ is positive or negative is the same as the question of how higher derivative corrections affect the masses of Kaluza-Klein harmonics when we reduce on S 2 . In particular, in a theory with sufficient supersymmetry one might expect that δM ℓ ≥ 0 for all modes.
In summary, we have shown that higher derivative corrections to the geometry and linear higher derivative corrections to the scalar field equation of motion do not lead to a qualitative change in the behaviour of linear scalar field perturbations at the Poincaré horizon of AdS 2 × S 2 . The dominant ℓ = 0 Aretakis instability is protected by the assumed shift symmetry of the scalar field. Higher derivative corrections can lead to small changes in the exponents of the power-law behaviour exhibited by higher ℓ modes but, for small α/Q, these corrections are small and so the ℓ = 0 instability remains dominant. There is no sign of any breakdown of effective field theory. Why do the higher derivative corrections to the equation of motion not become large? The reason can be traced to the fact that these corrections appear only via n ψ in (2.23). This structure is a consequence of general covariance, i.e., the fact that the higher derivative terms do not depend on anything except the background geometry. The high degree of symmetry of the background geometry then greatly restricts the form of the higher derivative terms in the action. Note in particular that general covariance forbids the appearance in the action of higher derivative terms evaluated in some geometrically preferred basis, such as the basis (Lie transported w.r.t. V ) that is used to exhibit the instability.
Full black hole solution
We have just seen that the higher derivative corrections do not cause a problem during the Aretakis instability in the near-horizon geometry. However, as we have just argued, this may be a consequence of the high degree of symmetry of the near-horizon geometry. It is not obvious that this result will still hold if we consider the less symmetric extremal RN geometry. Furthermore, the above analysis did not incorporate nonlinear corrections to the equations of motion (except via correcting the background geometry). In this section we will address both of these deficiencies by considering higher derivative corrections during the Aretakis instability in the full extreme RN geometry.
We will assume that the extremal RN solution can be corrected to give a static, spherically symmetric, solution to all orders in α, with Φ = const. For a large black hole, i.e., one with α/Q ≪ 1, the effect of corrections to this background solution should be small so we will neglect them in this section. We will focus on the effect of the higher derivative corrections to the equations of motion during the Aretakis instability. For effective theory to remain valid, these terms should remain small, giving perturbative corrections to the 2-derivative theory.
If the higher derivative terms become larger than the 2-derivative terms then effective field theory breaks down. So in this section we will investigate whether or not this is the case.
We will consider all of the equations of motion, not just the scalar field equation of motion.
First we note that coupled gravitational and electromagnetic perturbations of the extreme RN black hole exhibit an Aretakis instability [7] but this is weaker than the massless scalar field instability in the sense that it requires more derivatives to see it. So we will continue to focus on the Aretakis instability driven by a massless scalar field. This instability is strongest in the spherically symmetric ℓ = 0 sector. So if higher derivatives are going to cause trouble it seems very likely that this will occur in the ℓ = 0 sector. Therefore we can simplify by restricting to spherical symmetry.
We recall the effect of nonlinearities in the 2-derivative theory. As discussed in the Introduction, the nonlinear evolution of the spherically symmetric instability in the 2-derivative theory was studied in Ref. [14], where it was shown that the initial perturbation can be finetuned so that the metric "settles down" to extreme RN on and outside the event horizon, with the scalar field on the horizon exhibiting the Aretakis instability. In other words, the "most unstable" behaviour exhibited by the nonlinear 2-derivative theory is to give a spacetime which, at late time, looks like a linear scalar field on a fixed extreme RN background. Motivated by these results, our strategy in this section will be to consider a spherically symmetric scalar field evolving in a fixed extreme RN background. We will perform a consistency check on the smallness of the higher derivative corrections to the equations of motion. To do this we will take the known results for the late time behaviour of the scalar field along the horizon in the 2-derivative theory, and use this to estimate the size of higher derivative corrections to the equation of motion. In particular, we can compare the size of the higher derivative terms to (possibly nonlinear) terms present in the 2-derivative theory.
In order for effective field theory to remain valid, the higher derivative terms must remain small compared to the 2-derivative terms.
The extremal Reissner-Nordstrom solution is a type D solution, i.e., the Weyl tensor has two pairs of coincident principal null directions, which are also principal null directions of the Maxwell field. It is convenient to employ the Geroch-Held-Penrose (GHP) formalism [20], which is well suited to situations in which one has a pair of preferred null directions.
This formalism is based on a null tetrad and enables all calculations to be reduced to the manipulation of scalar quantities. In the metric ( In the GHP formalism, there is a freedom to change the basis (2.35) so that the two null directions are preserved. One possibility is to rescale the null vectors (referred to as a boost) where λ is a real function. The other is to rotate the spatial basis vectors (referred to as a spin) m → e iθ m ;m → e −iθm . The GHP formalism is designed to maintain convariance under boosts and spins. A privileged role is played by objects which transform covariantly, i.e., objects with definite boost and spin weight. Not all connection components transform covariantly. Those that do take the following values in the extreme RN background: The GHP scalars ρ, ρ ′ have boost weights 1, −1 respectively, and both have zero spin.
Since the background spacetime is type D, the only non-zero components of the Weyl tensor and Maxwell field are those with vanishing boost and spin weights The non-vanishing Ricci tensor components have boost weight zero and are determined by φ 1 . The GHP formalism introduces derivative operators with definite spin/boost weights. In the extreme RN background, they are given by where η is a GHP scalar with boost weight b and spin s, and ǫ, γ and β are Newman-Penrose spin coefficients. The operators þ, þ ′ have zero spin and carry boost weight 1, −1 respectively, and the operators , ′ have zero boost weight and carry spin 1, −1 respectively.
Finally we will need to use commutators of these derivative operators. Acting on a quantity of boost weight b and spin s, in the extreme RN background these are given by (2.42) Now we return to considering the higher-derivative corrected equations of motion in the extreme RN spacetime with a dynamical spherically symmetric scalar field. Consider a boost-weight B component of one of the equations of motion. We will determine the vdependence of higher derivative corrections to this component on the horizon at late time.
In the GHP formalism, all quantities are written as scalars so any higher-derivative term can be written in the form XZ where X is constructed entirely from the background GHP scalars and their derivatives, and Z is constructed entirely from the scalar field and its derivatives. We can write Z = Z 1 . . . Z N where each Z i consists of GHP derivatives acting on Φ. Spherical symmetry implies that none of these derivatives can be or ′ . To see this, note that any Z i can be written asD 1 . . .D p D 1 . . . D q Φ, or the corresponding expression with replaced by ′ , whereD i ∈ {þ, þ ′ , , ′ } and D i ∈ {þ, þ ′ }, for some p, q ≥ 0. But D 1 . . . D q Φ has spin 0, so, using spherical symmetry, it is annihilated by and ′ . Hence any Z i involving or ′ must vanish.
Next, using the commutator [þ, þ ′ ], we can order þ and þ ′ derivatives in Z i so that þ derivatives appears to the left of þ ′ derivatives. So there is no loss of generality in assuming that each Z i has the form þ j þ ′k Φ. Recall that we assumed that Φ is derivatively coupled but one might wonder whether commutators could generate terms without GHP derivatives. However this is not possible: [þ, þ ′ ] acting on derivatives of Φ gives a result involving derivatives of Φ whereas [þ, þ ′ ] acting on Φ gives zero (because Φ has zero boost weight). Hence commutators cannot give rise to terms involving Φ without derivatives so we must have j + k ≥ 1. Now on the horizon we have δ = 0 so we can replace þ with ∂ v in þ j þ ′k Φ and converting (2.16) to GHP notation gives where b = j − k is the boost weight of this term and ǫ ∈ {0, 1} with ǫ = 0 if k = 0 or k ≥ j + 1 and ǫ = 1 otherwise. Taking a product of N such terms gives ǫ i and we have used the fact that XZ has boost weight B, so we have where B X is the boost weight of X. Now, since X is constructed from background quantities, it is independent of v hence we have We will now show that if B X > 0 then X vanishes on the horizon. The scalar X can be written as X = X 1 . . . X M , where each X i consists of GHP derivatives acting on some GHP scalar ω associated to the background spacetime, i.e., ω ∈ {ρ, ρ ′ , Ψ 2 , φ 1 , φ * 1 }. Note that all of these quantities have zero spin and are spherically symmetric. This means that we can argue as above to show that or ′ derivatives cannot appear in X i . Using commutators, we can assume that X i has the form þ j þ ′k ω. Furthermore, since we can replace þ by ∂ v on the horizon, and the GHP scalars are all v-invariant, the expression þ j þ ′k ω vanishes when evaluated on the horizon unless j = 0. So any X i that is non-vanishing on the horizon must be of the form þ ′k ω. This has boost weight b ω − k where b ω is the boost weight of ω. Note that the possible ω all have non-positive boost weight, with the exception of ω = ρ. So if ω is anything except ρ then X i , if non-vanishing on the horizon, must have non-positive boost weight. If ω is ρ then b ω = 1 but, since ρ vanishes on the horizon, we need k ≥ 1 to construct a non-vanishing expression. Hence X i also has non-positive boost weight in this case. Therefore we have proved that all X i that are non-vanishing on the horizon must have non-positive boost weight. This proves that if X is non-vanishing on the horizon then B X ≤ 0.
Let's apply this to the Einstein equation, which has components with |B| ≤ 2. (Note that spherical symmetry implies that the B = ±1 components are trivial.) In the 2-derivative theory, the RHS of the Einstein equation involves the energy-momentum tensor of the scalar field. We'll denote this 2-derivative energy momentum tensor as T Φ µν . Equation B X > N + E − 2. But we've just seen that non-vanishing X on the horizon requires B X ≤ 0 so we'd need N < 2−E for our higher derivative term to dominate. However, we've assumed that all terms in the action are at least quadratic in the scalar field, which implies that all terms in the Einstein equation have N ≥ 2 (or N = 0 but the latter don't depend on the scalar field and hence don't depend on v). Hence it is not possible for higher derivatives to become large compared to the 2-derivative terms in the Einstein equation. The "worst" that can happen is that the higher derivative terms exhibit the same scaling with v as the 2-derivative terms. This happens when N = 2, E = 0 and B X = 0. Such terms scale in the same way as the 2-derivative terms but they will be suppressed by powers of the small quantity α/Q.
The same argument can be applied to the scalar field equation of motion, which has B = 0. A typical 2-derivative term in this equation of motion is þþ ′ Φ ∼ v −2 . So for a higher derivative term to dominate we would need B X − N − E > −2 i.e., B X > N + E − 2 so again we'd need N < 2 − E for consistency with B X ≤ 0. Our assumption that the scalar field appears at least quadratically in the action implies that N ≥ 1 in the scalar field equation of motion. There is now a non-trivial solution to these inequalities given by N = 1, E = 0 and B X = 0. However, such terms are excluded by our assumption of a shift symmetry. To see this, note that with N = 1, Z is linear in the scalar field, i.e., of the form þ j þ ′k Φ and with B = B X = 0 this term must have boost weight j − k = 0 so j = k. Now E = 0 implies ǫ = 0 which is only possible if j = k = 0, i.e., there are no derivatives acting on Φ. However we explained above that such a term is forbidden by our assumption that the scalar field has a shift symmetry. So in fact the "worst" terms are ones for which the higher derivative terms exhibit the same v −2 scaling as the two-derivative terms but are suppressed by powers of α/Q. Such terms can have either N = 1 or N = 2. With N = 1 these terms have Z of the form þΦ or þ j þ ′j Φ with j ≥ 1. With N = 2 these terms have Z of the form Since components of the Maxwell equation have |B| ≤ 1 we see that these terms decay at late time along the horizon.
These calculations demonstrate that there is no obvious failure of effective field theory on the horizon at late time. Although certain higher derivatives of the scalar field become large on the event horizon at late time, this does not imply that higher derivative corrections to the equation of motion become large compared to the 2-derivative terms. This is because, in the equations of motion, the "bad" derivatives are always multiplied by "good" terms which are decaying, or by terms X which vanish on the horizon. The reason for this can be traced back to general covariance. This implies that the quantities X appearing in the higher derivative terms are constructed only from GHP scalars associated to the background solution. In particular X depends only on the background fields and not on any additional structure such as a preferred basis. So, just as we found for the near-horizon geometry, it is general covariance which prevents a breakdown of effective field theory.
Extremal Kerr
In this section we will discuss the scalar field instability at the horizon of an extremal Kerr black hole, first discussed by Aretakis in the axisymmetric case and extended to the non-axisymmetric case in Ref. [9]. Our goal is to understand whether higher derivative corrections could become important during this instability. As for extremal RN, we will start by analyzing this in the near-horizon geometry before turning to the full black hole solution.
Near-horizon analysis
As explained above, the near-horizon AdS 2 × S 2 geometry of an extremal RN black hole provides a simplified setting in which to study the Aretakis instability [8]. Here we will consider the near-horizon extremal Kerr (NHEK) geometry [21] as a simplified setting to study the Aretakis instability of extremal Kerr. In fact our main motivation here is to go beyond the (axisymmetric) Aretakis instability and consider non-axisymmetric perturbations of extremal Kerr, as discussed in Ref. [9].
In the axisymmetric case, the results of Ref. [9] do not see the dominant Aretakis instability, behaving as in (1.1). This is because the approach of Ref. [9] cannot incorporate the presence, in the initial data, of outgoing radiation at the event horizon, so all the Aretakis constants are zero. Under such circumstances there is still an instability but it requires an extra derivative to see it [6], and this "subleading" instability was reproduced in Ref. [9]. For non-axisymmetric perturbations, Ref. [9] found an instability stronger than that discovered by Aretakis, with the first derivative of the scalar field generically growing along the horizon. However, since the approach of Ref. [9] cannot model outgoing radiation initially present at the event horizon one might wonder whether the inclusion of such radiation would make the non-axisymmetric instability even worse. This is something that we can investigate using the methods of this section. We will assume that the extremal Kerr solution with M ≫ α can be corrected to all orders in α to give an extremal black hole solution of the theory (2.2) and that this corrected solution has vanishing Maxwell field and constant scalar field. The general results of Ref. [22] imply that the near-horizon geometry of this black hole has SL(2, R) × U(1) symmetry and the metric can be written as an S 2 fibred over AdS 2 : where k(α) is a constant and Λ i are smooth functions on the sphere parameterized by (θ, ϕ).
The coordinates {T, R, ϕ} are then the near horizon descendants of the time, radial and axial coordinates of extreme Kerr in Boyer-Lindquist form. For nonzero α, we will refer to (3.1) as the α-NHEK geometry.
The coordinates {T, R, θ, ϕ} cover a patch of α-NHEK which is analogous to the Poincaré patch in AdS 2 . We can covert to global coordinates (described in appendix A) to obtain what we will call the global α-NHEK geometry. The AdS 2 part of this geometry is depicted by the infinite vertical strip in figure 1. One of the SL(2, R) generators of the isometry group can be taken to be the translations in global time τ (see appendix A), that is-shifts up and down the 'global α-NHEK' strip in figure 1. We will make use below of a translation with ∆τ = π/2 which in Poincaré corresponds to the transformation (see also [24], [25]) 3) is an isometry: the metric in the new coordinates is precisely of the same form as (3.1), replacing {T, R, ϕ} → {t, r, χ}. We will start by considering the wave equation in the above geometry, i.e. we neglect higher derivative corrections to the scalar equation of motion in this section. Supposing initial data for ψ is specified on some surface in the near-horizon region, for example T − 1/R = const. < 0 as seen in figure 1, we would like to study the resulting solution.
Ref. [23] studied perturbations of near-horizon geometries of the α-NHEK type, and in particular it was shown that they are separable and the wave equation reduces to the equation of a massive charged scalar in AdS 2 with a homogeneous electric field. To see this, use the ansatz and Fourier decompose along the φ direction as Define the effective AdS 2 metric and gauge field where∇ is the covariant derivative on AdS 2 and q = −mk is the effective electric charge. Then the equation governing X(T, R) is where λ is the eigenvalue of the angular equation where∇ is the covariant derivative on the transverse S 2 with metric defined by setting dT = dR = 0 in (3.1). The operator O can be shown to be self-adjoint w.r.t. an appropriate inner product so its eigenvalues are real and the eigenfunctions form a complete set on S 2 [23]. Hence there is no loss of generality in decomposing ψ as in (3.4). In general, these eigenfunctions can be labelled by a pair of integers (ℓ, m) with |m| ≤ ℓ just as for standard spherical harmonics.
Equation (3.8) describes a scalar field with charge q and squared mass µ 2 = λ+q 2 in AdS 2 with an electric field. The electric field is homogeneous because the corresponding Maxwell 2-form is proportional to the AdS 2 volume form. If one separates variables, i.e., assumes e −iωT time dependence then solutions of the radial equation have two possible behaviours as R → ∞, given by [21,26,27] As R → ∞, a general superposition of such modes will behave as for some functions f ± (T ). For well-defined dynamics we need to impose boundary conditions at R = ∞. If h is real then a natural choice is to impose "normalizable" boundary conditions, i.e., f + ≡ 0. In NHEK this is the case for axisymmetric modes, i.e., m = 0, for which λ = ℓ(ℓ + 1) and hence h = ℓ + 1 [21]. However, if λ < −1/4 then h is complex. For NHEK this occurs for non-axisymmetric modes with |m| ∼ ℓ. In this case it is not clear what boundary conditions should be imposed (see Refs. [21,26,27] for discussions of this issue).
We will assume that for complex h one can obtain well-posed dynamics with a boundary condition that fixes some linear relation between f + and f − .
Notice that the axisymmetric modes will have real h in α-NHEK. This is because the associated eigenvalues λ are non-negative in NHEK so small higher derivative corrections to the background geometry cannot push λ below −1/4 in α-NHEK. Hence the higher derivative corrections to the background geometry will lead to small real shifts in h. This will not happen for the ℓ = 0 mode, i.e., the constant mode on S 2 , which continues to have λ = 0 and h = 1 in α-NHEK. For the non-axisymmetric modes, it is possible that a NHEK mode with λ slightly larger than −1/4 (hence real h) might correspond to an α-NHEK mode with λ slightly less than −1/4 (hence complex h).
The idea now is that we can determine the late time behaviour of the scalar field along the Poincaré horizon in α-NHEK simply from a coordinate transformation. We consider the Poincaré horizon r = 0 in the coordinates (t, r, θ, χ). We shift to ingoing Eddington-Finkelstein coordinates (v, r, θ, so that the metric is now regular at the Poincaré horizon: Late time along the Poincaré horizon corresponds to r = 0, v → ∞. From Fig. 1, this can be seen to correspond to the limit R → ∞, T → 0 in the original coordinates. So we can determine the late-time behaviour of the scalar field by transforming (3.11) to the new coordinates. Doing this, including the angular dependence e imϕ S(θ), gives Here we have transformed to the new coordinates and taken the limit v → ∞ with rv fixed.
In figure 1, rv represents the angle of approach to the center of the dotted circle as the limit v → ∞ is taken. On the horizon we have rv = 0 but it is convenient to allow for non-zero rv because it enables us to see explicitly the r-dependence of ψ at late time near the horizon.
For the modes with real h, which includes the axisymmetric modes, we impose normalizable boundary conditions f + (0) = 0. From the above expression we have and where D denotes angular derivatives. 7 Note that when h is real we have h ≥ 1/2.
For modes with complex h, which are non-axisymmetric, we have h = 1/2 + iζ where ζ is real. We then have |ψ| horizon ∼ v −1/2 (3.17) and This is precisely the late time behaviour discovered for the full extremal Kerr solution in Ref. [9]. As mentioned above, the approach of Ref. [9] cannot incorporate the effects of outgoing radiation initially present at the event horizon (or non-vanishing Aretakis constants in the axisymmetric case) so one might wonder whether the presence of such radiation could change the results, perhaps leading to even slower decay. Our analysis allows for outgoing radiation initially present at the event horizon and our results agree with those of Ref. [9] when h is complex. This suggests that inclusion of the initial outgoing radiation does not lead to slower decay. Of course it would be desirable to confirm this using an analysis in the full black hole spacetime rather than just the near-horizon geometry. The analysis of this section could also be generalised to fields of higher spin, where one would need to supplement the transformation (3.3) with a tetrad rotation (c.f. [28]).
Linear higher-derivative corrections in near-horizon geometry
So far we have studied a massless scalar in the α-NHEK geometry, i.e., we have incorporated higher derivative corrections to the background geometry but not to the scalar equation of motion. In this section we will investigate the effects of the linear higher derivative corrections to the massless scalar equation of motion. We cannot consider nonlinear corrections to the equations of motion because it is known that 2-derivative nonlinearities (i.e. backreaction) tend to destroy the NHEK asymptotics [26,27].
We will proceed as we did for AdS 2 × S 2 in section 2.3, i.e, expanding the action to quadratic order in ψ, substituting in the expansion of ψ in terms of spheroidal harmonics on S 2 : and then integrating over S 2 to obtain an action governing the charged fields X λm in AdS 2 with a homogeneous electric field as in (3.6). The axisymmetry of the background implies that modes corresponding to harmonics with different values of m will decouple from each other in the action. However, the θ-dependence of the background will lead to coupling of the modes with different values of λ (but the same m) in the dimensional reduction of the higher derivative terms. Because of the SL(2, R) symmetry of the background, the resulting action for the fields of charge q = −km will have the form (integrating by parts so derivatives act on X and notX) where g 2 is the AdS 2 metric in (3.6) and (since the action is real) Our assumption that ψ is derivatively coupled implies that X 00 cannot appear without derivatives in the above action. This is because Y 00 is constant and hence eliminated by angular derivatives, so X 00 must be acted on by AdS 2 derivatives. Therefore we must have c 0λ00 = 0 and hence c 00λ ′ 0 = 0. It is convenient to define a vector X m with components X λm and Hermitian matrices C mn with components c mλλ ′ n . The action can then be written Since C mn is the coefficient of a term with 2n derivatives we must have 8 for some dimensionless HermitianC mn . For n = 1, 0 we can use the known equation of motion in the 2-derivative theory and the fact that the higher derivative corrections start at O(α 2 ) to deduce and that where J m has components j mλλ ′ = − λ + (mk) 2 δ λλ ′ (3.26) In the above we are ignoring a possible overall factor in the action. We now repeat the strategy of section 2.3 using a field redefinition to eliminate the higher derivative terms in S m . Henceforth we suppress the m index and write where D n are dimensionless matrices depending on α/M. Substituting this into the action whereX is a vector with componentsX λ and E n are Hermitian matrices. The first few of these are We now want to choose the unknown matrices D n so that E n vanishes for n ≥ 2. This can be done order by order in α/M. We start with E 2 = 0 which, using (3.24), gives . Plugging this back into E 2 = 0 then determines the O(α 2 /M 2 ) part of D 2 . Repeating this process order by order we achieve E n = 0 for all n ≥ 2. The action has become (3.31) E 1 is Hermitian so we can diagonalize it with a unitary matrix U: where K is real and diagonal. Furthermore we have and we can choose U = I + O(α 2 /M 2 ). Since K is positive definite we can write K = L † L for a positive definite real diagonal matrix L = I + O(α 2 /M 2 ). We now bring the kinetic term to canonical form with a final field redefinition: where we have defined the Hermitian "mass matrix" where we have reinstated the m indices.
We have now included the effects of higher derivative terms both via the correction to the background geometry, and via the correction to the linearized equation of motion for the scalar field. Both effects can be incorporated simply by a perturbative shift λ → λ + O(α 2 /M 2 ) in the value of λ that appears in the effective AdS 2 equation of motion. This translates into a perturbative shift of the conformal weights (3.10) which determine the decay rates at late time along the Poincaré horizon.
Recall that the slowest decaying modes are non-axisymmetric with complex h, i.e., λ < −1/4. For these modes, a small perturbative shift in λ will still result in complex h and hence the decay results (3.17) and (3.18) will still hold. So we conclude that higher derivative corrections to the background and linear higher derivative corrections to the scalar equation of motion do not change the rate of decay of the slowest decaying NHEK modes.
For modes with real h, the shift in λ will result in a small correction to the decay rates (3.15), (3.16), similar to what happens to the ℓ > 0 modes in AdS 2 × S 2 , as described in section 2.3. However (after field redefinitions) the λ = 0, m = 0 mode does not suffer a correction, as a consequence of the shift symmetry of the scalar field. To see this, note that X 00 does not appear in the "mass" term in (3.31) because of c 00λ ′ 0 = c 0λ00 = 0. Hence varying (3.31) w.r.t.X 00 gives an equation of motion (E 1 ) 0λ D 2X 0λ = 0. So (E 1 ) 0λX0λ = X 00 + O(α 2 /M 2 ) satisfies a decoupled equation of motion with λ = 0.
In summary, our near-horizon analysis, taking into account all higher derivative corrections to the background, and linear higher derivative corrections to the equation of motion, indicates that higher derivative corrections do not make the scalar field instability of Ref. [9] any worse. So the near-horizon analysis does not indicate any breakdown of effective field theory at late time at the horizon. As for AdS 2 × S 2 , the reason for this is that general covariance combined with the SL(2, R) symmetry greatly restricts the possible form of the higher derivative terms in the action (3.20).
Higher derivative corrections in full black hole geometry
We have shown that higher derivative corrections do not cause a problem during the scalar field instability in the NHEK geometry. However, this may be a consequence of the high symmetry of this near-horizon geometry. It is not obvious that this result will still hold if we consider the less symmetric extremal Kerr geometry. Furthermore, the above analysis did not incorporate nonlinear corrections to the equations of motion (except via correcting the background geometry). In this section we will address both of these deficiencies by considering higher derivative corrections during the scalar field instability in the full extremal Kerr geometry.
We will perform calculations analogous to the calculations we performed for extremal Reissner-Nordstrom in section 2.4. We will assume that the extreme Kerr solution can be corrected, to all orders in α, to give a stationary, axisymmetric, neutral BH solution. Assuming that the BH is large, α ≪ M, will allow us to neglect the corrections to the background in this section's analysis. We will then take the known behaviour of a massless scalar field on the horizon of an extremal Kerr black hole and use it to compare the size of higher derivative corrections to the equation of motion to the size of two-derivative terms. 9 There is an immediate problem with this investigation. In the two derivative Einsteinscalar theory, there has been no study of backreaction of the scalar field instability of extremal Kerr. So if the effects of two derivative nonlinearities are not understood, how are we to understand higher derivative terms? In this section we will simply assume, in analogy with the extremal RN case, that the "worst" behavior in the nonlinear two-derivative theory is that the spacetime settles down to extremal Kerr on and outside the event horizon, with the scalar field behaving just like a linear field in the extreme Kerr spacetime. With this assumption, we will determine the behaviour of higher derivative terms in the equations of motion.
We start with the Kerr metric written in ingoing Kerr coordinates (v, r, θ,χ): The event horizon is at r = M i.e. δ = 0. We now convert to co-rotating coordinates (v, r, θ, χ) defined byχ In these coordinates, ∂/∂v is tangent to the horizon generators. The Kerr solution is type D and we choose a null tetrad based on the two repeated principal null directions. In coordinates (v, r, θ, χ), the basis is The GHP connection scalars are: The type D property means that the only non-vanishing GHP curvature scalar is The GHP derivative operators are given by Commutators of these derivatives acting on a quantity of boost weight b and spin s are given by Consider a component of the equations of motion which has boost weight B. As in section 2.4 we note that any higher derivative term has the form XZ where X is constructed from background GHP quantities and Z is constructed from the scalar field and its derivatives.
We write Z = Z 1 . . . Z N where each Z i consists of GHP derivatives acting on Φ. Using GHP commutators we can arrange these derivative so that Z i has the form þ j þ ′k l ′m Φ. The assumed shift symmetry implies that, before using commutators, Φ always appears with derivatives acting on it. From the explicit form of the commutators, we see that a commutator acting on derivatives of Φ gives terms involving derivatives of Φ and a commutator acting on Φ also gives derivatives of Φ (because Φ has b = s = 0). Hence commutators cannot generate terms involving Φ without derivatives so j + k + l + m ≥ 1. We assume that Φ is composed of all possible harmonics in extreme Kerr, so the late time behaviour is dominated by the non-axisymmetric modes with m ∼ ℓ, i.e. the modes with complex h, for which, on the horizon at late time [9] |∂ j v ∂ k r D l Φ| horizon ∼ v k−j−1/2 (3.45) where D denotes angular derivatives. Since þ ∼ ∂ v on the horizon, this implies that þ j þ ′k l ′m Φ ∼ v k−j−1/2 = v −b−1/2 where b = j − k is the boost weight of this term. From this we have where B X is the boost weight of X. Since X is constructed from background quantities, it is independent of v so we also have XZ| horizon ∼ v B X −B−N/2 (3.47) We will now show that if B X > 0 then X vanishes on the horizon. We write X = X 1 . . . X M where each X i consists of GHP derivatives acting on some GHP scalar ω associated to the background spacetime, i.e., ω ∈ {τ, τ ′ , ρ, ρ ′ , Ψ 2 } (or complex conjugates of these). Using commutators we can assume that X i has the form þ j þ ′k l ′m ω. Since þ ∼ ∂ v on the horizon, and all GHP scalars are v-invariant, it follows that this expression vanishes on the horizon unless j = 0. So any X i that is non-vanishing on the horizon must have the form þ ′k l ′m ω, which has boost weight b ω − k where b ω is the boost weight of ω. Note that b ω ≤ 0 unless ω = ρ. So if ω = ρ then X i , if non-vanishing on the horizon, must have non-positive boost weight. If ω = ρ then b ω = 1 but, since ρ vanishes on the horizon, we need k ≥ 1 to construct a non-vanishing expression. So in this case, X i also has non-positive boost weight if non-vanishing on the horizon. It follows that B X ≤ 0 if X is non-vanishing on the horizon. Now let's apply this to the Einstein equation, which has components with |B| ≤ 2. In the 2-derivative theory, the energy momentum tensor of Φ has components which scale as v −B−1 at late time along the horizon. So in order for the higher derivative term (3.47) to become large compared to this 2-derivative term we would need B X − B − N/2 > −B − 1, i.e., 2B X > N − 2. But non-vanishing X require B X ≤ 0 so this is possible only if N < 2, which contradicts our assumption that the scalar field appears at least quadratically in the action and hence quadratically in the Einstein equation. So it is not possible for the higherderivative terms to become large compared to the 2-derivative terms. The worst that can happen is for the higher derivative terms to exhibit the same v-dependence as the 2-derivative terms, suppressed by powers of the small quantity α/M. This happens if N = 2 and B X = 0.
For the scalar field equation of motion we have B = 0 and typical 2-derivative terms are þþ ′ Φ ∼ ′ Φ ∼ v −1/2 . So for a higher derivative term to become large compared to this we would need B X − N/2 > −1/2, i.e., 2B X > N − 1. But B X ≤ 0 and in the scalar field equation of motion we have N ≥ 1 so this is not possible. The worst that can happen is when N = 1 and B X = 0, i.e., linear, boost weight zero, higher derivative corrections with Z of the form þ j þ ′j l ′m Φ. These exhibit the same late time v-dependence as the 2-derivative terms but they are suppressed by powers of α/M. In summary, our conclusions are the same as for extremal RN. Even though the nonaxisymmetric scalar field instability of extremal Kerr is worse than the axisymmetric Aretakis instability, we have found that, at the horizon, higher derivative corrections remain small compared to 2-derivative terms. Once again the underlying reason for this can be traced to general covariance, which greatly restricts the form of the higher derivative terms.
Specifically, it implies that the quantity X in the above argument is constructed from GHP scalars associated to the background geometry. This gave us the restriction B X ≤ 0 which eliminates dangerous higher derivative terms in the above argument. We should emphasize that the analysis of this section started from the assumption that, when we include backreaction in the 2-derivative theory, the "worst" than can happen is that the spacetime "settles down" to extremal Kerr, with the scalar field evolving at late time as a test field in the extremal Kerr background. If this assumption is incorrect then our analysis would no longer apply. So clearly the most important issue here is to understand this backreaction in the 2-derivative theory. | 15,131.2 | 2017-09-27T00:00:00.000 | [
"Physics"
] |
A Graphene-Based THz Wave Duplexer and Filter: Switching via Gate Biasing
: This work introduces a Graphene-based multi-layer reconfigurable device as a wave duplexer in the THz frequency range. Adjusting transmitting and reflecting parts of incident waves alongside controlling absorption provides the interesting capability to select target waves in different frequencies. The proposed device includes periodic graphene patterns on both sides of silicon dioxide as substrate. Additionally, the patterns are biased differently compared to conventional patterns which makes it possible to achieve two distinct behaviors versus frequency. Exploiting equivalent circuit models (ECM) for graphene and dielectric, the whole device is modeled by passive RLC circuits. According to simulation results, the proposed device can transmit and reflect incident THz waves at desired frequencies in 0.1 THz to 30 THz which makes it an ideal candidate for manipulating THz waves in terms of transmission and reflection.
provides carrier frequencies in telecommunication protocols much higher than those used in current communication systems (wireless and commercial). It has also been used to increase data transmission capacity for short-range (short-distance) communications. The various applications of the terahertz band have made this area one of the most interesting research topics [4][5][6][7][8]. Implementing such applications requires materials to realize terahertz devices. Researchers interested in the terahertz band have studied the properties of materials with different dimensions. In terms of common three-dimensional materials, the electrical conductivity and conductivity of silica and gold have been investigated in the terahertz band. In addition, a range of two-dimensional materials such as graphene, phosphorene, and silicon have also been investigated [9].
Among these two-dimensional materials, graphene has been the subject of research in various fields due to its unique mechanical, thermal, and electromagnetic properties. One of the attractive properties of graphene as a semiconductor-metal material is its unique conductivity, which can be controlled by external voltage bias or static magnetic bias. Now, along with the terahertz frequency band and graphene material, we need a way to model and realize such devices. To help design accurate structure, it is essential to provide numerical modeling methods that take into account the properties of graphene [9], [30][31][32][33][34].
One of the simplification approaches to solving electromagnetic problems is to describe the elements as circuit elements. In the meantime, relatively accurate equations are presented to describe the graphene layer.
In 2014 and 2015, references [10] and [11] presented efficient circuit models for graphene nano-strips and graphene nano-disks, respectively, which consider the effects of physical parameters such as the geometry and electron relaxation time along with the effects of bias voltage has taken into account.
According to many previous works, the accuracy of the circuit model performance has been compared with numerical methods, most of which report a very good agreement with the slight error of the circuit model method versus numerical methods [10], [11]. Based on this, it can be concluded that the circuit model approach can be a substitute for time-consuming and complex numerical methods. In this regard, three graphene patterns are focused.
Nano-strips, nano-disks, and graphene continuous plates are modeled in series with resistor-inductor-capacitor circuit elements. This modeling, along with knowing the dielectric circuit model, leads to the calculation of the structure impedance. Matching this impedance with the impedance of the outer space of the structure means the transfer of maximum power from the external environment to the structure (maximum power in an absorber means complete absorption in the structure). In other words, the impedance adaptation is equivalent to the full absorption at the frequency at which this adaptation occurred. Thus, this approach paves the way for the design of graphene-based devices by simplifying an electromagnetic phenomenon to the level of an impedance matching problem [18][19][20][21][22][23].
In this way, here two graphene patterns on both sides of dielectric form a reconfigurable device that can select the desired reflection and/or transmission waves. Section 2 describes the proposed device in detail with corresponding circuit representation. Also, section 3 provides simulations results while section 4 concludes the work as the conclusion.
Proposed Device:
The proposed device is illustrated in Fig. 1, includes two periodic arrays of Graphene on both sides of a dielectric. Each Graphene pattern consists of a triple-bias scheme. According to [14], bias equalities can force patterns to experience different periods, and consequently, a degree of freedom is provided to adjust device reactions in desired frequencies. In addition, based on [12][13][14][15], [35][36][37], circuit representation of graphene patterns is reported with excellent accuracy compared to full numerical simulations. According to Fig. 1, the incident wave divides into three parts: Transmitted, Absorbed, and Reflected parts. By minimizing the absorption part, interesting functions can be realized via transmission and reflection parts. In this way, the proposed structure can act similar to a duplexer which can transmit or reflect target waves.
, , P, W, hr and D are dielectric thickness, disk radius, disk period, ribbons width, and ribbons period respectively. It should be noted that the illustration is symbolic while all disks are the same in shape and all ribbons have equal width. Also, based on [15], periods of disks and ribbons can vary be corresponding to bias equalities. According to [10] and [14], each pattern can be modeled via a resistor, inductor, and capacitor. These elements are related to some physical constants and geometrical sizes as bellow: (1) To obtain n q , and 1n q , physical parameter such as , P, r W , and D , must be designed and then referred to Table 2. Z impedance sees to graphene disks at the end of the line. Therefore, 1 Z calculated by Eq.
(3). Then Eq. (4) is obtained based on the transmision line concept [16]. Also, the dielectric can be modeled via (5). And the input impedance is calculated as below: where the definition of parameters of Eq. (3) to Eq. (6) are reported in Table 3: Table 3. Definition of parameters.
Definition Parameter
The wave propagation constant in the dielectric substrate. Based on Eq. (6), the input impedance of the proposed structure is calculable versus frequency. All parameters in the proposed structure have their effects on the impedance. So finding the impedance sensitivity to parameters is interesting to obtain distinct responses. (7) In this way, with an iterative algorithm, two sets of chemical potentials are obtained as Mode A and Mode B which are tabulated in the next section.
Simulation Results:
According to Eq. (8), changing gate bias leads to changing chemical potential which forces the device to react differently [17]. In this way, two operational modes are defined with corresponding chemical potentials, tabulated in Table 4. Table 1.
The described device is simulated via two sets of bias values. The corresponding chemical potentials for each bias set are tabulated in Table 4. Table 4. Two sets of bias values for the proposed device.
First operational mode (eV) Second operational mode (eV) Similarly, the second mode of operation is reported in Fig. 6 and Fig. 7.
According to simulation results, both operational modes are in an acceptable condition against geometrical variations while the device can be switched easily between two modes just by setting bias values. Finally, the comparison table is reported in Table 5.
Conclusion:
Using triple-bias graphene patterns of ribbons and disks, a reconfigurable device is presented which can manipulate THz radiation in terms of transmission, reflection, and absorption. Changing external gate bias stimulates the device to shift operational region against frequency. The device is modeled by circuit model elements to simplify simulation.
The sensitivity of the proposed device versus geometrical parameters is reported which verifies the superior performance of the design. Such a reconfigurable device is in great demand for several functions in optical systems. It can be used as a wave duplexer with the capability of tuning via gate biasing.
Data Availability Statement:
The data that support the findings of this study are available from the corresponding author upon reasonable request. | 1,836.8 | 2021-12-01T00:00:00.000 | [
"Physics",
"Engineering"
] |
Lipocalin 2 as a Putative Modulator of Local Inflammatory Processes in the Spinal Cord and Component of Organ Cross talk After Spinal Cord Injury
Lipocalin 2 (LCN2), an immunomodulator, regulates various cellular processes such as iron transport and defense against bacterial infection. Under pathological conditions, LCN2 promotes neuroinflammation via the recruitment and activation of immune cells and glia, particularly microglia and astrocytes. Although it seems to have a negative influence on the functional outcome in spinal cord injury (SCI), the extent of its involvement in SCI and the underlying mechanisms are not yet fully known. In this study, using a SCI contusion mouse model, we first investigated the expression pattern of Lcn2 in different parts of the CNS (spinal cord and brain) and in the liver and its concentration in blood serum. Interestingly, we could note a significant increase in LCN2 throughout the whole spinal cord, in the brain, liver, and blood serum. This demonstrates the diversity of its possible sites of action in SCI. Furthermore, genetic deficiency of Lcn2 (Lcn2−/−) significantly reduced certain aspects of gliosis in the SCI-mice. Taken together, our studies provide first valuable hints, suggesting that LCN2 is involved in the local and systemic effects post SCI, and might modulate the impairment of different peripheral organs after injury.
Introduction
Spinal cord injury (SCI) is a devastating event that causes life-long health restrictions including paralysis, loss of sensation and vegetative functions, pain, and psychological impairment [1]. Despite many efforts, there is presently no comprehensive treatment protocol available to effectively treat this injury, mainly owed to the complexity of nerve fiber tract destructions, neuronal death, and poor restoration capacities of physiological function [2,3]. In SCI, the primary injury refers to the initial physical damage of the spinal cord (SC), which is accompanied by hemorrhage, ischemia, and local neuronal death, while the secondary injury phase is characterized by progressive damage of the SC, demyelination, astrogliosis, and neuroinflammation [4][5][6][7][8][9]. Progressing neuroinflammation, which is a major hallmark of the secondary injury, is mainly initiated through activation of astrocytes and microglia, which are key cells in the maintenance of homeostasis in the CNS, and further boosted and perpetuated by infiltrated neutrophils and macrophages [10][11][12][13].
The activation of astrocytes and microglia, so-called astrogliosis and microgliosis, respectively, influences the disease outcome in SCI on various levels [14,15]. Astrocytes are the predominant subtype of glial cells in the CNS. Under physiological conditions, they protect neurons through the uptake of excessive neurotransmitters, i.e., glutamate, maintain the integrity of the blood-brain barrier and participate in synaptic stability, plasticity, and reorganization [16,17]. When being activated, astrocytes become hypertrophic and develop extended processes [6]. Reactive astrocytes are a central component of the glial scar which is formed around the injury site in the secondary injury phase [18]. Glial scar formation affects the healing process and can remain chronically for up to several decades in patients who suffered from SCI [14]. The glial scar limits the spread of inflammation but, at the same time, impedes axonal regeneration [19][20][21]. Under pathological conditions such as traumatic SCI, reactive astrocytes promote cytotoxic edema formation and ischemia through an upregulation of aquaporin 4 [16]. Furthermore, they are an integral component of local immune responses by producing and secreting a wide range of cytokines and chemokines [22,23]. It has been shown that the phenotype of reactive astrocytes varies, and it has been assumed that astrocytes can differentiate in the direction of either a more pro-inflammatory A1 or a more anti-inflammatory A2 polarization state [24]. A1 polarized astrocytes express pro-inflammatory cytokines and contribute to neuronal death, whereas A2 polarized astrocytes stimulate CNS recovery and repair [24,25].
The neuroinflammation in SCI is regulated by expression of pro-inflammatory and anti-inflammatory cytokines, chemokines, and other mediators, which are mainly synthesized by glial cells. The glycoprotein lipocalin 2 (LCN2) is considered a key mediator of immune responses in general and particularly in neurodegenerative diseases [26][27][28][29][30][31]. It has been shown that LCN2, which is upregulated at the lesion site of the SC, is produced by astrocytes after SCI [26]. Furthermore, Lcn2-deficient mice reveal better functional outcomes, a lower expression of chemokines, and a reduced extent of secondary injury after SCI in comparison to wild-type mice [26]. In general terms, LCN2 plays an important role in iron transport and homeostasis and promotes the defense against bacterial infections [27,32]. Furthermore, it has been demonstrated in vitro that LCN2 has toxic effects on neurons and regulates the expression of pro-inflammatory cytokines and chemokines [30,33]. It has further been stated that LCN2 promotes the shifting of the polarization of microglia and astrocytes towards proinflammatory phenotypes in vitro [25,34].
It has been shown that SCI causes pathological processes in various parts of the body, which were not directly affected by the injury. In patients suffering from SCI cognitive dysfunction, inflammation-associated neurodegeneration of brain tissue and an impaired functional brain recovery are commonly observed [35][36][37]. In addition to the interaction between different parts of the CNS, a further important issue related to neural injury is the communication with peripheral organ systems, i.e., "CNS-organ cross talk". There are preliminary findings which show that in SCI, a defined communication axis exists between SC and liver suggesting that the liver might exhibit mechanisms that influence neuroinflammation in the SC [38,39].
Due to the limited treatment options in SCI, it is important to identify new possible drug targets. As we suggest LCN2 to influence SCI pathology, it is of interest to examine its effects on astrocytes, which play a central role in SCI pathology. Additionally, we wanted to get a first impression of whether LCN2 might participate in the systemic effects of SCI. In the present study, we have analyzed the time course of local Lcn2 expression post SCI and its influence on the activation and polarization of astrocytes. Furthermore, since LCN2 is also secreted in a paracrine and endocrine fashion, we analyzed the amount of LCN2 in blood and other peripheral organs post SCI.
The mice were housed and handled in accordance with the guidelines of the Federation for European Laboratory Animal Science Associations (FELASA) under standard laboratory conditions. The procedures were approved by the Review Board for the Care of Animal Subjects of the district government (North Rhine-Westphalia, Germany) and performed according to international guidelines on the use of laboratory mice (Az 81-02.04.2018.A227). Lcn2 −/− mice of the 7-day group were approved by the Review Board for the Care of Animal Subjects of the district government (ethic No. 962055, Tehran, Iran).
The WT mice were obtained from Janvier Labs (Saint-Berthevin Cedex, France); information on the genetic identification of the WT mice is available on their homepage (https:// www. janvi er-labs. com/ en/ an-optim al-manag ementof-genet ics-a-unique-conce pt/). The Janvier Labs' colony belongs to a genetically tested and characterized founding pair (genetic analysis 640,000 SNPs) that is identical to that of the C57BL/6JRj. Lcn2-deficient mice (Lcn2 −/− ), which have already been used in other studies from our group, were kindly provided by Tak W Mak (University of Toronto, Canada) and colleagues [30,40,41]. In these mice, a targeted mutation has been introduced to disrupt the Lcn2 coding region, including exons 1-5, with a PGK-neo cassette, thus leading to a functional knockout in all tissues including the CNS. For breeding, pairs of homozygous mice were used.
Spinal Cord Injury
General anesthesia was initiated with isoflurane (2-3 vol%) in an anesthetic chamber. During surgery, isoflurane (1.5-2 vol%) was further administered via a face mask. Intraoperative analgesia was attained through injection of buprenorphine (0.05-1 mg/kg s.c.) 30 min preoperatively. After the exposure of the spinal column (T7-T10), a laminectomy of T8 was performed. A standardized injury of the SC at this level was induced by contusion (Infinite Horizons Spinal Cord Impactor) with a force of 60 kdyn. After inducing the SCI, the surgical site was sutured in layers and the mice were injected subcutaneously with sterile saline. Postoperative care involved the daily manual emptying of the bladder until spontaneous urination returned. In the control group, after preoperative analgesia and general anesthesia as described above, the spinal column (T7-T10) was exposed, and a laminectomy of T8 was performed. The surgical site was then sutured in layers, and the mice were injected subcutaneously with sterile saline. Through this approach, we aim to preclude possible falsifications of the results caused by the mere surgical procedure.
BBB Scoring
To assess functional recovery and locomotion deficits after SCI, the mice were scored in an open field according to Basso, Beattie, and Bresnahan (BBB) locomotion rating scale of 0 (complete paralysis) to 21 (normal) as previously described [42]. The scale assesses hind limb movements, body weight support, forelimb to hind limb coordination, and whole-body movements.
Tissue Preparation
At defined time points after SCI (6,12,24,72, h and 7 days), the mice were transcardially perfused with ice-cold PBS for molecular biological and protein biochemical studies. The sham-operated mice, which served as control, were finalized after 24 h. The whole SC was prepared and divided into three parts, in the following referred to as rostral, central (lesion site), and caudal region. The three spinal cord regions were separated at the level of T3 and L1. This ensures that the caudal and rostral regions are located at a sufficient distance of several millimeters from the visually visible lesion area.
In addition, motor and sensory cortex and left liver lobe were prepared. The tissues were immediately snap frozen in liquid nitrogen and kept at − 80 °C until further processing.
Molecular Biological Analysis
For RNA isolation, the tissues were placed in homogenization tubes containing 1.4-mm beads. Samples were homogenized at 5000 × g for 15 s. RNA was isolated by phenol-chloroform extraction using peqGold RNA Tri-Fast (PeqLab, Erlangen, Germany). Total RNA amount and purity are determined using 260/280 ratios of optical densities (Nanodrop 1000, PeqLab, Erlangen, Germany). cDNA was obtained by reverse transcription using M-MLV reverse transcription (RT) kit and random hexanucleotide primers (Invitrogen, Carlsbad, USA). Gene expression levels were analyzed with real-time reverse transcription-PCR (Bio-Rad, Feldkirchen, Germany) using SensiMix™ SYBR® & Fluorescein Kit (Meridian Bioscience, Cincinnati, USA). Distilled water was used instead of cDNA as negative control. Primer sequences and individual annealing temperatures are shown in Table 1. Results were evaluated using Bio-Rad CFX manager (Bio-Rad, Feldkirchen, Germany) and were normalized to cyclophilin A and Hsp90 as reference genes.
The target gene expression was calculated using the ΔΔCt method [43].
Protein Biochemical Analysis
Sampled tissues were mechanically disrupted in RIPA buffer (pH 8.0) supplemented with a protease inhibitor cocktail (Complete Mini, Roche Diagnostics, Grenzach-Wyhlen, Germany). Protein concentrations were determined using the PierceTM BCA Protein Assay kit (Thermo Fisher Scientific, Waltham, USA) according to the manufacturer's protocol. Per sample, a total of 20 μg protein was separated in a 14% SDS polyacrylamide gel by gel electrophoresis and transferred to a PVDF (polyvinylidene difluoride) membrane. The blots were blocked in 5% milk in tris-buffered saline (TBS, pH 7.4) and then incubated overnight (at 4 °C) in primary antibodies rabbit anti-LCN2 in 5% milk and rabbit anti-GAPDH in 5% milk (used antibodies are listed in Table 2). An appropriate secondary antibody (goat anti-rabbit IgG (H + L)-HRP) was applied for 2 h (RT). Signals were analyzed via chemiluminescence detection (Westar Supernova, XLS 3,0100, Cyanagen, Bologna, Italy), visualized (Fusion Solo X, Vilber, Eberhardzell, Germany) and subjected to densitometry analysis using Image J. Results were normalized to GAPDH as reference protein.
Immunohistochemistry
For immunohistochemistry (IHC), 5-µm thick sections of SC, brain, and liver were rehydrated, and antigens were unmasked by heating in Tris/EDTA (pH 9.0) buffer for 20 min. After blocking with 5% normal goat serum in PBS, the sections were incubated overnight (4 °C) with rabbit anti-LCN2, or rabbit anti-ALDH1L1 respectively, diluted in 5% normal serum in PBS. Slides were incubated for 30 min in 0.3% H 2 O 2 (in PBS) followed by incubation with goat anti-rabbit IgG (H&L) diluted in 5% normal serum in PBS for 1 h (RT). Afterwards, an incubation with ABC-solution (both parts 1:50, VECTASTAIN Elite ABC Kit (Standard), Vector Labs, Burlingame, USA) diluted in PBS for 1 h (RT) followed. For double immunofluorescence labeling, sections were blocked with IFF buffer, containing BSA, FCS and 1 × PBS, for 1 h and incubated overnight (4 °C) with rabbit anti-LCN2 diluted in IFF buffer. The slides were incubated with donkey anti-rabbit 488 diluted in IFF buffer for 1 h (RT) followed by an incubation with goat anti-GFAP, respectively mouse anti-IBA1, rat anti-CD44, or rat anti-CD31 diluted in IFF buffer overnight (4 °C). Finally, the slides were incubated with donkey anti-goat 594, respectively donkey anti-mouse 594 or goat anti-rat 555 in IFF buffer for 1 h (RT). As negative controls, slices of the examined tissue, which were not incubated with the respective primary antibodies, were used. Apart from that, the negative controls were treated like the stained slices.
Statistical Analysis
A total of 59 WT animals were used for the experiments containing 43 animals for qPCR analysis. Twenty-four out of the 43 animals were also used for western blot analysis. Samples from 39 animals were subjected to ELISA. For immunohistochemistry staining, we used slices from 16 animals. A total of 20 Lcn2 −/− mice were used for qPCR analysis and 4 for immunohistochemistry. Per group, 4 animals and 3 sections per animal at a distance of 100 µm were stained.
GraphPad Prism 8 (GraphPad Software Inc., San Diego, USA) was used for statistical analysis. Brown-Forsythe test was performed to test for equal variances and normal distribution was tested with Shapiro-Wilk's test. If necessary, data were transformed via Boxcox for homoscedasticity. One-way ANOVA followed by Dunnett's post hoc test or two-way ANOVA followed by Tukey's post hoc test was used for parametric data. Non-parametric data (Lcn2 mRNA in sensory and motor cortex and LCN2 concentration in blood serum) were analyzed with the Kruskal-Wallis test followed by Dunn's multiple comparisons. WT and Lcn2 −/− data from BBB scoring were compared by two-way ANOVA with Geisser-Greenhouse correction. All data are given as arithmetic means ± standard errors of the mean (SEM). The p values were set as *p ≤ 0.05, **p ≤ 0.01, ***p ≤ 0.001, and ****p ≤ 0.0001, respectively #p ≤ 0.05, ##p ≤ 0.01, ###p ≤ 0.001, and ####p ≤ 0.0001.
Results
In a first set of experiments, we aimed at investigating whether traumatic SCI leads to an increase in Lcn2 expression within the SC and other peripheral organs. Figure 1 shows a significant and stepwise increase of LCN2 in the central region (injury site) of the SC. mRNA expression immediately rose within the first 6 h post injury reaching maximum level at 24 h post SCI and then rapidly declined at 7 days (Fig. 1a). LCN2 protein levels, which were examined by western blot, revealed a similar time course and profile with a short delay compared to mRNA expression, peaking at 72 h post SCI (Fig. 1b/c). To investigate the distribution and localization of LCN2 positive cells in injured SC, immunostaining against LCN2 was performed. Immunohistochemistry showed high numbers of LCN2-positive cells, especially in the gray matter, in the central lesion region 24 h post injury compared to that of the control group (Fig. 1d/e/f). Double immunofluorescence staining revealed that LCN2 signals are associated with GFAPpositive astrocytes occasionally (Fig. 1j/k). IBA1-positive microglia (Fig. 1g) did not co-localize with LCN2 staining in the SC. By staining against the leukocyte marker CD44 and LCN2, co-expression of CD44 could be seen in most LCN2-positive cells 24 h post SCI (Fig. 1h/i).
To determine whether Lcn2 is upregulated in rostral and caudal parts of the spinal cord, we measured the Lcn2 mRNA levels in these regions. The results indicate a massive upregulation of Lcn2 during the initial 6 h in the rostral part which persists until 72 h (Fig. 2a). In the caudal area, we observed a steady upregulation of Lcn2 during the first 7 days post SCI (Fig. 2b).
In a next step, we have examined whether Lcn2 is upregulated in the brain post SCI. In both examined brain regions, sensory and motor cortex, significantly elevated Lcn2 mRNA levels were already present 6 h post SCI and thereafter declined (Fig. 2c/d). Protein levels were analyzed in the sensory cortex (Fig. 2e/f) showing a similar expression pattern as Lcn2 mRNA. As shown in Fig. 2, our results from immunofluorescence staining against LCN2 revealed no reactivity in the brain slices of sham-operated mice, but LCN2 + cells occurred, mainly around vessels, after SCI (Fig. 2g/h). To identify these cells as endothelial cells, we performed immunofluorescence double staining, which showed a clear co-localization of LCN2 with the endothelial marker CD31 (Fig. 2i).
Furthermore, we analyzed the LCN2 concentration in blood serum via ELISA, which was significantly elevated around 12 and 24 h post SCI, reaching a ~ 19-fold increase at its peak (Fig. 3a). In addition, we assessed a potential Lcn2 upregulation in the liver. Here, Lcn2 mRNA (Fig. 3b) and protein (Fig. 3c/d) were significantly elevated from 6 h post SCI on, reaching a peak at 12 h and decreasing again from then on. After immunofluorescence staining, almost no LCN2-immunoreactive cells could be seen in the control group, whereas scattered immunoreaction was detectable after SCI (Fig. 3e/f).
A common phenomenon after SCI is astrogliosis. Since astrocytes are one of the LCN2-producing cell types, we aimed to correlate the expression of astroglial markers (Gfap, vimentin, serpina3n) and Lcn2 in the central SC region (Supplementary Fig. 1a-c). Like Lcn2 mRNA, Gfap, vimentin, and serpina3n mRNA show a significant and progressive increase from 6 h post SCI on. Serpina3n, like Lcn2, reaches its peak at 24 h and decreases from then on, whereas Gfap and vimentin levels proceed to rise.
In order to understand the influence of LCN2 on the pathological scenario after SCI better, we have included animals with a general Lcn2 deficiency (Lcn2 −/− ) in our study. To assess locomotor impairment and recovery of WT and Lcn2 −/− mice after SCI, we used BBB scoring (Fig. 4a). Control animals of both genotypes were all rated with a score of 21, demonstrating their unimpaired condition. At 24 and 72 h, no significant differences in BBB scoring could be seen between the two genotypes. After 7 days, Lcn2 −/− mice reached a mean score of ~ 8, indicating sweeping with no weight support or plantar placement of the paw with no weight support, whereas the mean score of ~ 4 in WT mice stands for only slight movement of all three joints of the hind limbs. The significantly higher scores of Lcn2 −/− mice at 7 days indicate better locomotor recovery compared to that of WT mice.
By comparing results from WT and Lcn2 −/− tissues, we could first demonstrate that the gene expression of the astrogliosis marker Gfap was reduced in the central SC region of Lcn2 −/− animals compared to that of WT mice at all examined time points, with a significant difference between the two genotypes at 24 h (Fig. 4b). Results from staining against the astrocyte marker ALDH1L1 (Supplementary Fig. 2) show fewer positive cells in the lesion region of Lcn2 −/− compared to WT at 7 days. Activated astrocytes can differentiate in the direction of a more pro-or a more anti-inflammatory state and consequently have varying effects on disease pathology. The functional polarization of astrocytes is well-acknowledged. Since LCN2 was shown to influence this polarization in vitro, we also addressed the influence of LCN2 on astrocyte polarization in SCI. Complement component 3 (C3) and sphingosine kinase 1 (Sphk1) were selected as representative markers for A1, respectively A2 astrocyte polarization states. Both markers are well-known and recognized and have been described in publications of renowned journals [24,44]. Therefore, we analyzed the gene expression profiles of these markers in the injured SC (Fig. 4c-f). In Fig. 4c, a value < 1 of the C3/Sphk1 (A1/A2) quotient indicates a prevalence of A2 during the first 24 h after SCI in the central region in WT mice. From 72 h on, values > 1 demonstrate a prevalence of A1. The underlying individual evaluation of C3 and Sphk1 in WT is not shown. These findings suggest that there are changes in the polarization of astrocytes after SCI.
In order to explore whether LCN2 influences, additionally to the extent of astrogliosis, also the functional polarization of astrocytes, we assessed the mRNA expression of the A1 and A2 markers stated above in Lcn2 −/− mice. In the central region of the SC, a significant decrease in A1 and A2 marker mRNA could be seen at 24 h and 72 h (A1), respectively at all examined time points (A2) in Lcn2 −/− mice (Fig. 4d/e). The A1/A2 (C3/Sphk1) quotient in Lcn2 −/− shows the same pattern as in WT mice with an initial decrease at 24 h followed by a subsequent increase (Fig. 4f).
To assess the effects of LCN2 on apoptosis rates, the ratio of Bax mRNA, an apoptotic marker, and Bcl2 mRNA, an anti-apoptotic marker, was evaluated in WT and Lcn2 −/− mice. As we expected, we observed a significant increase of the Bax/Bcl2 quotient, indicating a pro-apoptotic state, in the central part of the SC at 24 and 72 h in WT mice (Fig. 4g). In contrast, the Bax/Bcl2 quotient did not change significantly compared to the control in the rostral and caudal region ( Supplementary Fig. 1d/e). In Lcn2 −/− mice, we observed only a slight reduction of Bax/Bcl2 ratios in the central SC region in comparison to WT mice, which did not reach a significant level (Fig. 4g).
Discussion
In the present study, we used a well-established SCI contusion mouse model to provide evidence that LCN2 is upregulated after SCI throughout the whole SC and not only in the primarily injured region. Beyond SC, we observed a LCN2induction in the cerebral cortex at both protein and mRNA levels. Interestingly, we show a marked increase of LCN2 in systemic circulation and also in liver in the early phase post SCI. Various studies have found a correlation between increased LCN2 levels and CNS disorders, such as multiple sclerosis and stroke [28,30,45,46]. Therefore, using Lcn2 −/− mice, we investigate the effect of Lcn2 deficiency on astrogliosis as a hallmark of SCI. Since the results show a significant reduction of Gfap, a decrease of astrogliosis in Lcn2-deficient mice might be concluded. Post SCI, astrocytes proliferate and undergo morphological changes which include hypertrophy and the development of extended processes [6,47]. Through the release of neurotrophic factors, astrocytes support neurons in SC and thus, impaired astrocytic function has major consequences for neuronal function [17,48]. In brain injury, the ablation of reactive astrocytes was found to lead to substantial neuronal degeneration [17]. Moreover, astrocytes limit the spread of inflammation after SCI, since they are one of the dominant cell types of the glial scar which forms after injury [21,47]. Furthermore, activated astrocytes can express a variety of cytokines, chemokines, and the respective receptors, and therefore play a pivotal role in the neuroinflammatory processes in SCI [47,49]. Furthermore, axonal regeneration is inhibited by the glial scar and chondroitin sulfate proteoglycans which are produced by reactive glial cells, including astrocytes [19,21]. In addition, these proteoglycans impede process outgrowth of oligodendrocytes and thereby disturb remyelination [50,51]. Based on the dual character of astrocytes, it has been suggested that they can be classified into a neurotoxic A1 and neuroprotective A2 phenotype [24,25]. Different factors, such as chemokines and cytokines, e.g., IL-1β, TNF-α, and IL-10, have been found to control the development of astrocytes in the direction of either phenotype [25,52,53]. One of the regulators of astrocyte polarization is LCN2 which supports the pro-inflammatory A1 phenotype and decreases the polarization in the direction of A2 in vitro by inhibiting IL-4-STAT6 signaling [25]. The influence of LCN2 on astrocyte polarization, morphology, and migration is an important aspect of its regulatory function in neuroinflammation [27,54]. LCN2 is involved in various pathological processes, such as stroke, metabolic inflammation, diabetes, and nonalcoholic steatohepatitis [30,31,55,56]. It promotes inflammation through induction of pro-inflammatory cytokines via release of high mobility group box 1, which binds to toll-like receptor 4 and induces oxidative stress by activation of NOX-2 signaling [55]. Furthermore, beyond its effect on activation and polarization of microglia, LCN2 supports the recruitment of inflammatory cells by the induction of CXCL10 secretion and release of the neutrophil-recruitment signal IL-8 [31,34,[57][58][59].
In the present study, we could demonstrate that SCI induces an increase of Lcn2 expression throughout the whole SC. As the cellular source of LCN2 in the CNS, previous studies have identified astrocytes and endothelial cells [26,60]. As already described in a stroke mouse model and in a SCI contusion mouse model similar to the animal model used in our study, we also recognized a co-localization of GFAP and LCN2 in some cells [26,30].
After SCI, leukocytes are recruited to the lesion region within hours [61]. According to our results from co-staining against CD44 and LCN2, the major source of local LCN2 seems to be infiltrated leukocytes, as most LCN2-positive cells are also positive for the leukocyte marker CD44 [62]. However, we could not prove the production of LCN2 by microglia in our animal model [63]. The triggers of LCN2 production in this context are, besides others, cytokines such as IL-6 and NF-kappa B activation [60,64].
Since LCN2 is secreted, and elevated concentrations can be found in the blood circulation under pathological conditions, like multiple sclerosis, intestinal inflammation, and arthritic diseases, it has been described as a biomarker in several pathologies [65]. In the present study, we show that the LCN2 concentration is significantly increased in the serum as a direct consequence of SCI, which might suggest this molecule as a potential biomarker for traumatic SCI. Furthermore, circulating LCN2 could be considered a part of the systemic inflammatory response (SIR) which affects the homeostasis of peripheral organs such as the liver, kidney, lung, and intestine. Thereby, it contributes to the pathogenesis of multiple organ dysfunction after SCI and supports secondary injury to the SC [38,[66][67][68][69]. In addition, we were able to detect elevated LCN2 levels in the brain and liver. This can have at least two reasons: As a first possibility, LCN2 might be produced in the respective tissue. This is supported by the fact that we have found significantly increased Lcn2 mRNA in both brain and liver. Additionally, the identification of LCN2 + cells in both tissues after IHC staining indicates a production of LCN2 by the resident cells. In the brain, we could identify endothelial cells as a cellular source of LCN2 by double immunofluorescence staining. One of the possible triggers of LCN2 production in the brain is cytokines. For example, the i.p. application of IL-6 induces LCN2 production by vascular cells in the brain in mice [60]. In adipocytes also, TNF-α and IL-1β trigger LCN2 production in vitro [70]. Since various cytokines have been shown to be upregulated in the blood stream after SCI, they might lead to an increase in LCN2 production in endothelial cells [71].
In the liver, hepatocytes and neutrophil granulocytes have been identified as cellular sources of LCN2 [72,73]. It has been demonstrated in vitro that the cytokine IL-1β induces LCN2 production in a NF-kappa B-dependent manner in both cell types [74][75][76]. Due to the structure of the hepatic tissue, hepatocytes and recruited neutrophils come into close contact with cytokines, reaching the liver via the hepatic artery which might induce LCN2 production [77]. Since we have found elevated LCN2 levels in serum post SCI, LCN2 might also, besides its production by resident cells, reach the brain and the liver via the bloodstream.
In the brain, LCN2 has different beneficial as well as harmful effects [78]. In the ischemic brain, LCN2 contributes to neuronal cell death by promoting neuroinflammation [79]. However, in an experimental model of multiple sclerosis, Lcn2-deficient mice exhibited increased disease severity, suggesting a neuroprotective role of LCN2 [46]. In liver pathology, the effects of LCN2 have been discussed controversially. In phases of acute liver injury, LCN2 plays an essential role in liver homeostasis and lipid metabolism and protects hepatocytes, whereas it promotes liver injury and hepatic steatosis in a model of alcoholic steatohepatitis [80][81][82].
In our studies, the decrease of the astrogliosis marker Gfap in Lcn2 −/− mice is a first, valuable hint at a possible promotion of astrogliosis by LCN2 in SCI [83]. This assumption is further supported by the decrease in astrogliosis in the central spinal cord region in Lcn2 −/− compared to WT at 7 days demonstrated by immunohistochemical staining shown in Supplementary Fig. 2. In vitro, it has already been demonstrated that Gfap expression is promoted by LCN2 [84]. However, according to our results, LCN2 does not affect the regulatory mechanism underlying the phenotypic polarization of activated astrocytes in our animal model. The promotion of the classical inflammatory activation of astrocytes by LCN2 has, up to now, been only confirmed in vitro and in an animal model of transient middle cerebral artery occlusion [25,85]. Eventually, the effect of LCN2 on astrocyte polarization depends on the underlying pathology.
So far, we base our conclusions regarding the influence of LCN2 on astrogliosis and astrocyte polarization on qPCR studies. Therefore, possible posttranslational modifications cannot be taken into account. This limitation has to be addressed in further studies. Nevertheless, we confirm a general positive effect of Lcn2 deficiency on the functional outcome in SCI based on BBB locomotor scoring. It is assumed that the elevated level of LCN2 after SCI may exacerbate axonal degeneration and contribute to poor neurological outcome by enhancing inflammatory cell infiltration and promoting neuronal apoptosis [26].
In summary, we found that SCI promotes the LCN2upregulation in SC, brain, blood circulation, and peripheral organs such as the liver. Consequently, LCN2 might play a role in systemic effects and multiple organ dysfunction in SCI pathology. The precise effect of LCN2 on peripheral organs has to be examined thoroughly to understand the resulting SCI-induced impairment of these tissues. As a local consequence of SCI pathology, LCN2 promotes specific aspects of astrogliosis, which suggests that LCN2 can be therapeutically targeted to modulate the reaction of astrocytes in certain pathologies such as SCI. Further studies are needed to elucidate the precise mechanisms responsible for astrocyte activation and polarization to better understand the role played by LCN2 in this process. | 7,028 | 2021-08-21T00:00:00.000 | [
"Biology",
"Environmental Science",
"Medicine"
] |