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451781 | https://en.wikipedia.org/wiki/Polyacrylonitrile | Polyacrylonitrile | Polyacrylonitrile (PAN) is a synthetic, semicrystalline organic polymer resin, with the linear formula (CH2CHCN)n. Almost all PAN resins are copolymers with acrylonitrile as the main monomer. PAN is used to produce large variety of products including ultra filtration membranes, hollow fibers for reverse osmosis, fibers for textiles, and oxidized PAN fibers. PAN fibers are the chemical precursor of very high-quality carbon fiber. PAN is first thermally oxidized in air at 230 °C to form an oxidized PAN fiber and then carbonized above 1000 °C in inert atmosphere to make carbon fibers found in a variety of both high-tech and common daily applications such as civil and military aircraft primary and secondary structures, missiles, solid propellant rocket motors, pressure vessels, fishing rods, tennis rackets and bicycle frames. It is a component repeat unit in several important copolymers, such as styrene-acrylonitrile (SAN) and acrylonitrile butadiene styrene (ABS) plastic.
History
Polyacrylonitrile (PAN) was first synthesized in 1930 by Hans Fikentscher and Claus Heuck in the Ludwigshafen works of the German chemical conglomerate IG Farben. However, as PAN is non-fusible, and did not dissolve in any of the industrial solvents being used at the time, further research into the material was halted. In 1931, Herbert Rein, head of polymer fiber chemistry at the Bitterfeld plant of IG Farben, obtained a sample of PAN while visiting the Ludwigshafen works. He found that pyridinium benzylchloride, an ionic liquid, would dissolve PAN. He spun the first fibers based on PAN in 1938, using aqueous solutions of quaternary ammonium sodium thiocyanate and aluminum perchlorate for the production process and considered other solvents including DMF. However, commercial introduction was delayed due to the wartime stresses on infrastructure, inability to melt the polymer without degradation, and solvents to allow solution processing were not known yet.
The first mass production run of PAN fiber was in 1946 by American chemical conglomerate DuPont. The German intellectual property had been stolen in Operation Paperclip. The product, branded as Orlon, was based on a patent filed exactly seven days after a nearly identical German claim. In the German Democratic Republic (GDR), industrial polyacrylonitrile fibre production was started in 1956 at the VEB Film- und Chemiefaserwerk Agfa Wolfen due to the preliminary work of the "Wolcrylon" collective (:de:Max Duch, Herbert Lehnert et al.). Prior to this, the preconditions for the production of the raw materials had been created at the Buna Werke Schkopau (Polyacrylonitrile) and Leuna works (Dimethylformamide). In the same year, the collective was awarded the GDR's National Prize II Class for Science and Technology for its achievements.
Physical properties
Although it is thermoplastic, polyacrylonitrile does not melt under normal conditions. It degrades before melting. It melts above 300 °C if the heating rates are 50 degrees per minute or above.
Glass transition temperature is around 95 °C and fusion temperature is at 322 °C. PAN is soluble in polar solvents, such as dimethylformamide, dimethylacetamide, ethylene and propylene carbonates, and in aqueous solutions of sodium thiocyanate, zinc chloride or nitric acid. Solubility parameters: 26.09 MPa1/2 (25 °C) are 25.6 to 31.5 J1/2 cm−3/2. Dielectric constants: 5.5 (1 kHz, 25 °C), 4.2 (1 MHz, 25 °C).Can behave as branched as well as linear polymer.
Synthesis
Most commercial methods for the synthesis of PAN are based on free radical polymerization of acrylonitrile. In most of the cases, 10% amounts of other vinyl comonomers are also used (1–10%) along with AN depending on the final application. Comonomers include acrylic acid, acrylamide, allyl compounds, and sulfonated styrene. Anionic polymerization also can be used for synthesizing PAN. For textile applications, molecular weight in the range of 40,000 to 70,000 is used. For producing carbon fiber higher molecular weight is desired.
In the production of carbon fibers containing 600 tex (6k) PAN tow, the linear density of filaments is 0.12 tex and the filament diameter is 11.6 μm which produces a carbon fiber that has the filament strength of 417 kgf/mm2 and binder content of 38.6%. This data is demonstrated in the Indexes for Experimental Batches of PAN Precursor and Carbon Fibers Made from It table.
Applications
Homopolymers of polyacrylonitrile have been used as fibers in hot gas filtration systems, outdoor awnings, sails for yachts, and fiber-reinforced concrete. Copolymers containing polyacrylonitrile are often used as fibers to make knitted clothing like socks and sweaters, as well as outdoor products like tents and similar items. If the label of a piece of clothing says "acrylic", then it is made out of some copolymer of polyacrylonitrile. It was made into the spun fiber at DuPont in 1942 and marketed under the name of Orlon. Acrylonitrile is commonly employed as a comonomer with styrene, e.g. acrylonitrile, styrene and acrylate plastics. Labelling of items of clothing with acrylic (see acrylic fiber) means the polymer consists of at least 85% acrylonitrile as the monomer. A typical comonomer is vinyl acetate, which can be solution-spun readily to obtain fibers that soften enough to allow penetration by dyes. The advantages of the use of these acrylics are that they are low-cost compared to natural fiber, they offer better sunlight resistance and have superior resistance to attack by moths.
Acrylics modified with halogen-containing comonomers are classified as modacrylics, which by definition contain more than PAN percentages between 35-85%. Incorporation of halogen groups increases the flame resistance of the fiber, which makes modacrylics suitable for the use in sleepwear, tents and blankets. Some mattresses also use them to meet the flame resistance requirements in North America. However, the disadvantage of these products is that they are costly and can shrink after drying.
PAN absorbs many metal ions and aids the application of absorption materials. Polymers containing amidoxime groups can be used for the treatment of metals because of the polymers’ complex-forming capabilities with metal ions.
PAN has properties involving low density, thermal stability, high strength and modulus of elasticity. These unique properties have made PAN an essential polymer in high tech.
Its high tensile strength and tensile modulus are established by fiber sizing, coatings, production processes, and PAN's fiber chemistry. Its mechanical properties derived are important in composite structures for military and commercial aircraft.
Carbon fiber
Polyacrylonitrile is used as the precursor for 90% of carbon fiber production. Approximately 20–25% of Boeing and Airbus wide-body airframes are carbon fibers. However, applications are limited by PAN's high price of around $15/lb.
Glassy carbon
Glassy carbon, a common electrode material in electrochemistry, is created by heat-treating blocks of polyacrylonitrile under pressure at 1000 to 3000 °C over a period of several days. The process removes non-carbon atoms and creates a conjugated double bond structure with excellent conductivity.
Support polymer
Divinylbenzene-crosslinked polyacrylonitrile is a precursor to ion exchange resins. Hydrolysis converts the nitrile groups to carboxylic acids. Amberlite IRC86 is one commercial product. These weakly acidic resins have high affinities for divalent metal ions like Ca2+ and Mg2+.
| Physical sciences | Polymers | Chemistry |
451839 | https://en.wikipedia.org/wiki/MESSENGER | MESSENGER | MESSENGER was a NASA robotic space probe that orbited the planet Mercury between 2011 and 2015, studying Mercury's chemical composition, geology, and magnetic field. The name is a backronym for Mercury Surface, Space Environment, Geochemistry, and Ranging, and a reference to the messenger god Mercury from Roman mythology.
MESSENGER was launched aboard a Delta II rocket in August 2004. Its path involved a complex series of flybys – the spacecraft flew by Earth once, Venus twice, and Mercury itself three times, allowing it to decelerate relative to Mercury using minimal fuel. During its first flyby of Mercury in January 2008, MESSENGER became the second mission, after Mariner 10 in 1975, to reach Mercury.
MESSENGER entered orbit around Mercury on March 18, 2011, becoming the first spacecraft to do so. It successfully completed its primary mission in 2012. Following two mission extensions, the spacecraft used the last of its maneuvering propellant to deorbit, impacting the surface of Mercury on April 30, 2015.
Mission overview
MESSENGERs formal data collection mission began on April 4, 2011. The primary mission was completed on March 17, 2012, having collected close to 100,000 images. MESSENGER achieved 100% mapping of Mercury on March 6, 2013, and completed its first year-long extended mission on March 17, 2013. The probe's second extended mission lasted for over two years, but as its low orbit degraded, it required reboosts to avoid impact. It conducted its final reboost burns on October 24, 2014, and January 21, 2015, before crashing into Mercury on April 30, 2015.
During its stay in Mercury orbit, the probe's instruments yielded significant data, including a characterization of Mercury's magnetic field and the discovery of water ice at the planet's north pole, which had long been suspected on the basis of Earth-based radar data.
Mission background
Previous missions
In 1973, Mariner 10 was launched by NASA to make multiple flyby encounters of Venus and Mercury. Mariner 10 provided the first detailed data of Mercury, mapping 40–45% of the surface. Mariner 10's final flyby of Mercury occurred on March 16, 1975. No subsequent close-range observations of the planet would take place for more than 30 years.
Proposals for the mission
In 1998, a study detailed a proposed mission to send an orbiting spacecraft to Mercury, as the planet was at that point the least-explored of the inner planets. In the years following the Mariner 10 mission, subsequent mission proposals to revisit Mercury had appeared too costly, requiring large quantities of propellant and a heavy lift launch vehicle. Moreover, inserting a spacecraft into orbit around Mercury is difficult, because a probe approaching on a direct path from Earth would be accelerated by the Sun's gravity and pass Mercury far too quickly to orbit it. However, using a trajectory designed by Chen-wan Yen in 1985, the study showed it was possible to execute a Discovery-class mission by using multiple, consecutive gravity assist, 'swingby' maneuvers around Venus and Mercury, in combination with minor propulsive trajectory corrections, to gradually slow the spacecraft and thereby minimize propellant needs.
Objectives
The MESSENGER mission was designed to study the characteristics and environment of Mercury from orbit. The scientific objectives of the mission were:
to characterize the chemical composition of Mercury's surface.
to study the planet's geologic history.
to elucidate the nature of the global magnetic field (magnetosphere).
to determine the size and state of the core.
to determine the volatile inventory at the poles.
to study the nature of Mercury's exosphere.
Spacecraft design
The MESSENGER spacecraft was designed and built at the Johns Hopkins University Applied Physics Laboratory. Science operations were managed by Sean Solomon as principal investigator, and mission operations were also conducted at JHU/APL. The MESSENGER bus measured tall, wide, and deep. The bus was primarily constructed with four graphite fiber / cyanate ester composite panels that supported the propellant tanks, the large velocity adjust (LVA) thruster, attitude monitors and correction thrusters, the antennas, the instrument pallet, and a large ceramic-cloth sunshade, measuring tall and wide, for passive thermal control. At launch, the spacecraft weighed approximately with its full load of propellant. MESSENGER's total mission cost, including the cost of the spacecraft's construction, was estimated at under US$450 million.
Attitude control and propulsion
Main propulsion was provided by the 645 N, 317 sec. Isp bipropellant (hydrazine and nitrogen tetroxide) large velocity assist (LVA) thruster. The model used was the LEROS 1b, developed and manufactured at AMPAC‐ISP's Westcott works, in the United Kingdom. The spacecraft was designed to carry of propellant and helium pressurizer for the LVA.
Four monopropellant thrusters provided spacecraft steering during main thruster burns, and twelve monopropellant thrusters were used for attitude control. For precision attitude control, a reaction wheel attitude control system was also included. Information for attitude control was provided by star trackers, an inertial measurement unit and six Sun sensors.
Communications
The probe included two small deep space transponders for communications with the Deep Space Network and three kinds of antennas: a high gain phased array whose main beam could be electronically steered in one plane, a medium-gain "fan-beam" antenna and a low gain horn with a broad pattern. The high gain antenna was used as transmit-only at 8.4 GHz, the medium-gain and low gain antennas transmit at 8.4 GHz and receive at 7.2 GHz, and all three antennas operate with right-hand circularly polarized (RHCP) radiation. One of each of these antennas was mounted on the front of the probe facing the Sun, and one of each was mounted to the back of the probe facing away from the Sun.
Power
The space probe was powered by a two-panel gallium arsenide/germanium solar array providing an average of 450 watts while in Mercury orbit. Each panel was rotatable and included optical solar reflectors to balance the temperature of the array. Power was stored in a common-pressure-vessel, 23-ampere-hour nickel–hydrogen battery, with 11 vessels and two cells per vessel.
Computer and software
The spacecraft's onboard computer system was contained in an Integrated Electronics Module (IEM), a device that combined core avionics into a single box. The computer featured two radiation-hardened IBM RAD6000s, a 25 megahertz main processor, and a 10 MHz fault protection processor. For redundancy, the spacecraft carried a pair of identical IEMs. For data storage, the spacecraft carried two solid-state recorders able to store up to one gigabyte each. The IBM RAD6000 main processor collected, compressed, and stored data from MESSENGER's instruments for later playback to Earth.
MESSENGER used a software suite called SciBox to simulate its orbit and instruments, in order to "choreograph the complicated process of maximizing the scientific return from the mission and minimizing conflicts between instrument observations, while at the same time meeting all spacecraft constraints on pointing, data downlink rates, and onboard data storage capacity."
Scientific instruments
Mercury Dual Imaging System (MDIS)
Included two CCD cameras, a narrow-angle camera (NAC) and a wide-angle camera (WAC) mounted to a pivoting platform. The camera system provided a complete map of the surface of Mercury at a resolution of , and images of regions of geologic interest at . Color imaging was possible only with the narrow-band filter wheel attached to the wide-angle camera.
Objectives:
Flyby Phase:
Acquisition of near-global coverage at ≈.
Multispectral mapping at ≈.
Orbital Phase:
A nadir-looking monochrome global photomosaic at moderate solar incidence angles (55°–75°) and or better sampling resolution.
A 25°-off-nadir mosaic to complement the nadir-looking mosaic for global stereo mapping.
Completion of the multispectral mapping begun during the flybys.
High-resolution () image strips across features representative of major geologic units and structures.
Principal investigator: Scott Murchie / Johns Hopkins University
Gamma-Ray Spectrometer (GRS)
Measured gamma-ray emissions from the surface of Mercury to determine the planet's composition by detecting certain elements (oxygen, silicon, sulfur, iron, hydrogen, potassium, thorium, uranium) to a depth of 10 cm.
Objectives:
Provide surface abundances of major elements.
Provide surface abundances of Fe, Si, and K, infer alkali depletion from K abundances, and provide abundance limits on H (water ice) and S (if present) at the poles.
Map surface element abundances where possible, and otherwise provide surface-averaged abundances or establish upper limits.
Principal investigator: William Boynton / University of Arizona
Neutron Spectrometer (NS)
Determined the hydrogen mineral composition to a depth of 40 cm by detecting low-energy neutrons resulting from the collision of cosmic rays with the minerals.
Objectives:
Establish and map the abundance of hydrogen over most of the northern hemisphere of Mercury.
Investigate the possible presence of water ice within and near permanently shaded craters near the north pole.
Provide secondary evidence to aid in interpreting GRS measured gamma-ray line strengths in terms of elemental abundances.
Outline surface domains at the base of both northern and southern cusps of the magnetosphere where the solar wind can implant hydrogen in surface material.
Principal investigator: William Boynton / University of Arizona
X-Ray Spectrometer (XRS)
Mapped mineral composition within the top millimeter of the surface on Mercury by detecting X-ray spectral lines from magnesium, aluminum, sulphur, calcium, titanium, and iron, in the 1–10 keV range.
Objectives:
Determine the history of the formation of Mercury
Characterize the composition of surface elements by measuring the X-ray emissions induced by the incident solar flux.
Principal investigator: George Ho / APL
Magnetometer (MAG)
Measured the magnetic field around Mercury in detail to determine the strength and average position of the field.
Objectives:
Investigate the structure of Mercury's magnetic field and its interaction with the solar wind.
Characterize the geometry and time variability of the magnetospheric field.
Detect wave-particle interactions with the magnetosphere.
Observe magnetotail dynamics, including phenomena possibly analogous to substorms in the Earth's magnetosphere.
Characterize the magnetopause structure and dynamics.
Characterize field-aligned currents that link the planet with the magnetosphere.
Principal investigator: Mario Acuna / NASA Goddard Space Flight Center
Mercury Laser Altimeter (MLA)
Provided detailed information regarding the height of landforms on the surface of Mercury by detecting the light of an infrared laser as the light bounced off the surface.
Objectives:
Provide a high-precision topographic map of the high northern latitude regions.
Measure the long-wavelength topographic features at mid-to-low northern latitudes.
Determine topographic profiles across major geologic features in the northern hemisphere.
Detect and quantify the planet's forced physical librations by tracking the motion of large-scale topographic features as a function of time.
Measure the surface reflectivity of Mercury at the MLA operating wavelength of 1,064 nanometers.
Principal investigator: David Smith / GSFC
Mercury Atmospheric and Surface Composition Spectrometer (MASCS)
Determined the characteristics of the tenuous atmosphere surrounding Mercury by measuring ultraviolet light emissions, and ascertained the prevalence of iron and titanium minerals on the surface by measuring the reflectance of infrared light.
Objectives:
Characterize the composition, structure, and temporal behavior of the exosphere.
Investigate the processes that generate and maintain the exosphere.
Determine the relationship between exospheric and surface composition.
Search for polar deposits of volatile material, and determine how are the accumulation of these deposits are related to exospheric processes.
Principal investigator: William McClintock / University of Colorado
Energetic Particle and Plasma Spectrometer (EPPS)
Measured the charged particles in the magnetosphere around Mercury using an energetic particle spectrometer (EPS) and the charged particles that come from the surface using a fast imaging plasma spectrometer (FIPS).
Objectives:
Determine the structure of the planet's magnetic field.
Characterize exosphere neutrals and accelerated magnetospheric ions.
Determine the composition of the radar-reflective materials at Mercury's poles.
Determine the electrical properties of the crust/atmosphere/environment interface.
Determine characteristics of the dynamics of Mercury's magnetosphere and their relationships to external drivers and their internal conditions.
Measure interplanetary plasma properties in cruise and in Mercury vicinity.
Principal investigator: Barry Mauk / APL
Radio Science (RS)
Measured the gravity of Mercury and the state of the planetary core by utilizing the spacecraft's positioning data.
Objectives:
Determine the position of the spacecraft during both the cruise and orbital phases of the mission.
Observe gravitational perturbations from Mercury to investigate the spatial variations of density within the planet's interior, and a time-varying component in Mercury's gravity to quantify the amplitude of Mercury's libration.
Provide precise measurements of the range of the MESSENGER spacecraft to the surface of Mercury for determining proper altitude mapping with the MLA.
Principal investigator: David Smith / NASA Goddard Space Flight Center
Mission profile
Launch and trajectory
The MESSENGER probe was launched on August 3, 2004, at 06:15:56 UTC by NASA from Space Launch Complex 17B at the Cape Canaveral Air Force Station in Florida, aboard a Delta II 7925 launch vehicle. The complete burn sequence lasted 57 minutes bringing the spacecraft into a heliocentric orbit, with a final velocity of 10.68 km/s (6.64 miles/s) and sending the probe into a 7.9 billion-kilometer (4.9 billion mi) trajectory that took 6 years, 7 months and 16 days before its orbital insertion on March 18, 2011.
Traveling to Mercury and entering orbit requires an extremely large velocity change (see delta-v) because Mercury's orbit is deep in the Sun's gravity well. On a direct course from Earth to Mercury, a spacecraft is constantly accelerated as it falls toward the Sun, and will arrive at Mercury with a velocity too high to achieve orbit without excessive use of fuel. For planets with an atmosphere, such as Venus and Mars, spacecraft can minimize their fuel consumption upon arrival by using friction with the atmosphere to enter orbit (aerocapture), or can briefly fire their rocket engines to enter into orbit followed by a reduction of the orbit by aerobraking. However, the tenuous atmosphere of Mercury is far too thin for these maneuvers. Instead, MESSENGER extensively used gravity assist maneuvers at Earth, Venus, and Mercury to reduce the speed relative to Mercury, then used its large rocket engine to enter into an elliptical orbit around the planet. The multi-flyby process greatly reduced the amount of propellant necessary to slow the spacecraft, but at the cost of prolonging the trip by many years and to a total distance of 7.9 billion kilometers (4.9 billion miles).
Several planned thruster firings en route to Mercury were unnecessary, because these fine course adjustments were performed using solar radiation pressure acting on MESSENGER's solar panels. To further minimize the amount of necessary propellant, the spacecraft orbital insertion targeted a highly elliptical orbit around Mercury.
The elongated orbit had two other benefits: It allowed the spacecraft time to cool after the times it was between the hot surface of Mercury and the Sun, and also it allowed the spacecraft to measure the effects of solar wind and the magnetic fields of the planet at various distances while still allowing close-up measurements and photographs of the surface and exosphere.
The spacecraft was originally scheduled to launch during a 12-day window that beginning May 11, 2004. On March 26, 2004, NASA announced the launch would be moved to a later, 15-day launch window beginning July 30, 2004, to allow for further testing of the spacecraft. This change significantly altered the trajectory of the mission and delayed the arrival at Mercury by two years. The original plan called for three fly-by maneuvers past Venus, with Mercury orbit insertion scheduled for 2009. The trajectory was changed to include one Earth flyby, two Venus flybys, and three Mercury flybys before orbit insertion on March 18, 2011.
Earth flyby
MESSENGER performed an Earth flyby one year after launch, on August 2, 2005, with the closest approach at 19:13 UTC at an altitude of 2,347 kilometers (1,458 statute miles) over central Mongolia. On December 12, 2005, a 524-second-long burn (Deep-Space Maneuver or DSM-1) of the large thruster adjusted the trajectory for the upcoming Venus flyby by 316 m/s.
During the Earth flyby, the MESSENGER team imaged the Earth and Moon using MDIS and checked the status of several other instruments observing the atmospheric and surface compositions and testing the magnetosphere and determining that all instruments tested were working as expected. This calibration period was intended to ensure accurate interpretation of data when the spacecraft entered orbit around Mercury. Ensuring that the instruments functioned correctly at such an early stage in the mission allowed opportunity for multiple minor errors to be dealt with.
The Earth flyby was used to investigate the flyby anomaly, where some spacecraft have been observed to have trajectories that differ slightly from those predicted. However no anomaly was observed in MESSENGER's flyby.
Two Venus flybys
On October 24, 2006, at 08:34 UTC, MESSENGER encountered Venus at an altitude of . During the encounter, MESSENGER passed behind Venus and entered superior conjunction, a period when Earth was on the exact opposite side of the Solar System, with the Sun inhibiting radio contact. For this reason, no scientific observations were conducted during the flyby. Communication with the spacecraft was reestablished in late November and performed a deep space maneuver on December 12, to correct the trajectory to encounter Venus in a second flyby.
On June 5, 2007, at 23:08 UTC, MESSENGER performed a second flyby of Venus at an altitude of , for the greatest velocity reduction of the mission. During the encounter, all instruments were used to observe Venus and prepare for the following Mercury encounters. The encounter provided visible and near-infrared imaging data of the upper atmosphere of Venus. Ultraviolet and X-ray spectrometry of the upper atmosphere were also recorded, to characterize the composition. The ESA's Venus Express was also orbiting during the encounter, providing the first opportunity for simultaneous measurement of particle-and-field characteristics of the planet.
Three Mercury flybys
MESSENGER made a flyby of Mercury on January 14, 2008 (making its closest approach of 200 km above the surface of Mercury at 19:04:39 UTC), followed by a second flyby on October 6, 2008. MESSENGER executed a final flyby on September 29, 2009, further slowing down the spacecraft. Sometime during the closest approach of the last flyby, the spacecraft entered safe mode. Although this had no effect on the trajectory necessary for later orbit insertion, it resulted in the loss of science data and images that were planned for the outbound leg of the fly-by. The spacecraft had fully recovered by about seven hours later. One last deep space maneuver, DSM-5, was executed on November 24, 2009, at 22:45 UTC to provide the required velocity change for the scheduled Mercury orbit insertion on March 18, 2011, marking the beginning of the orbital mission.
Orbital insertion
The thruster maneuver to insert the probe into Mercury's orbit began at 00:45 UTC on March 18, 2011. The 0.9 km/s (0.5 mi./sec.) braking maneuver lasted about 15 minutes, with confirmation that the craft was in Mercury orbit received at 01:10 UTC on March 18 (9:10 PM, March 17 EDT). Mission lead engineer Eric Finnegan indicated that the spacecraft had achieved a near-perfect orbit.
MESSENGERs orbit was highly elliptical, taking it within of Mercury's surface and then away from it every twelve hours. This orbit was chosen to shield the probe from the heat radiated by Mercury's hot surface. Only a small portion of each orbit was at a low altitude, where the spacecraft was subjected to radiative heating from the hot side of the planet.
Primary science
After MESSENGER'''s orbital insertion, an eighteen-day commissioning phase took place. The supervising personnel switched on and tested the craft's science instruments to ensure they had completed the journey without damage. The commissioning phase "demonstrated that the spacecraft and payload [were] all operating nominally, notwithstanding Mercury's challenging environment."
The primary mission began as planned on April 4, 2011, with MESSENGER orbiting Mercury once every twelve hours for an intended duration of twelve Earth months, the equivalent of two solar days on Mercury. Principal Investigator Sean Solomon, then of the Carnegie Institution of Washington, said: "With the beginning today of the primary science phase of the mission, we will be making nearly continuous observations that will allow us to gain the first global perspective on the innermost planet. Moreover, as solar activity steadily increases, we will have a front-row seat on the most dynamic magnetosphere–atmosphere system in the Solar System."
On October 5, 2011, the scientific results obtained by MESSENGER during its first six terrestrial months in Mercury's orbit were presented in a series of papers at the European Planetary Science Congress in Nantes, France. Among the discoveries presented were the unexpectedly high concentrations of magnesium and calcium found on Mercury's nightside, and the fact that Mercury's magnetic field is offset far to the north of the planet's center.
Extended mission
In November 2011, NASA announced that the MESSENGER mission would be extended by one year, allowing the spacecraft to observe the 2012 solar maximum. Its extended mission began on March 17, 2012, and continued until March 17, 2013. Between April 16 and 20, 2012, MESSENGER carried out a series of thruster manoeuvres, placing it in an eight-hour orbit to conduct further scans of Mercury.
In November 2012, NASA reported that MESSENGER had discovered a possibility of both water ice and organic compounds in permanently shadowed craters in Mercury's north pole. In February 2013, NASA published the most detailed and accurate 3D map of Mercury to date, assembled from thousands of images taken by MESSENGER. MESSENGER completed its first extended mission on March 17, 2013, and its second lasted until April 2015. In November 2013, MESSENGER was among the numerous space assets that imaged Comet Encke (2P/Encke) and Comet ISON (C/2012 S1). As its orbit began to decay in early 2015, MESSENGER was able to take highly detailed close-up photographs of ice-filled craters and other landforms at Mercury's north pole. After the mission was completed, review of the radio ranging data provided the first measurement of the rate of mass loss from the Sun.
Discovery of water, organic compounds and volcanism
On July 3, 2008, the MESSENGER team announced that the probe had discovered large amounts of water present in Mercury's exosphere, which was an unexpected finding. In the later years of its mission, MESSENGER also provided visual evidence of past volcanic activity on the surface of Mercury, as well as evidence for a liquid iron planetary core. The probe also constructed the most detailed and accurate maps of Mercury to date, and furthermore discovered carbon-containing organic compounds and water ice inside permanently shadowed craters near the north pole.
Solar System portrait
On February 18, 2011, a portrait of the Solar System was published on the MESSENGER website. The mosaic contained 34 images, acquired by the MDIS instrument during November 2010. All the planets were visible with the exception of Uranus and Neptune, due to their vast distances from the Sun. The MESSENGER "family portrait" was intended to be complementary to the Voyager family portrait, which was acquired from the outer Solar System by Voyager 1 on February 14, 1990.
View of a total lunar eclipse
On October 8, 2014 from 9:18 UTC to 10:18 UTC, MESSENGER took 31 images, taken two minutes apart, of the Earth and the Moon, as the Moon underwent a total lunar eclipse. MESSENGER was 107 million kilometers (66 million miles) from the Earth at the time of the lunar eclipse. The Earth is about 5 pixels across and the Moon is just over 1 pixel across in the field of view of the NAC, with about 40 pixels distance between them. The images are zoomed by a factor of two and the Moon's brightness has been increased by a factor of about 25 to show its disappearance more clearly. This was the first observation of a lunar eclipse, of Earth's Moon, in history to be viewed from another planet.
End of mission
After running out of propellant for course adjustments, MESSENGER entered its expected terminal phase of orbital decay in late 2014. The spacecraft's operation was extended by several weeks by exploiting its remaining supply of helium gas, which was used to pressurize its propellant tanks, as reaction mass. MESSENGER continued studying Mercury during its decay period. The spacecraft crashed onto the surface of Mercury on April 30, 2015, at 3:26 p.m. EDT (19:26 GMT), at a velocity of , probably creating a crater in the planet's surface approximately wide. The spacecraft was estimated to have impacted at 54.4° N, 149.9° W on Suisei Planitia, near the crater Janáček. The crash occurred at a place not visible from Earth at the time, and thus was not detected by any observers or instruments. NASA confirmed the end of the MESSENGER mission at 3:40 p.m. EDT (19:40 GMT) after NASA's Deep Space Network did not detect the spacecraft's reemergence from behind Mercury.[[File:PIA19449-PlanetMercury-MESSENGER-Images-First-20110329-Last-20150430.jpg|thumb|center|600px|{{center|MESSENGERs first (March 29, 2011) and last (April 30, 2015) images from Mercury's orbit (impact details).}}]]
| Technology | Unmanned spacecraft | null |
17275317 | https://en.wikipedia.org/wiki/Anthoxanthin | Anthoxanthin | Anthoxanthins (flavones and flavonols) are a type of flavonoid pigments in plants. Anthoxanthins are water-soluble pigments which range in color from white or colorless to a creamy to yellow, often on petals of flowers. These pigments are generally whiter in an acid medium and yellowed in an alkaline medium. They are very susceptible to color changes with minerals and metal ions, similar to anthocyanins.
Uses
As with all flavonoids, anthoxanthins have antioxidant properties and are important for nutrition. They are sometimes used as food additives to add color or flavor to foods. One of the most well-known anthoxanthins is quercetin, which is found in many fruits and vegetables, including capers, red onions, and kale.
In addition to their use as food additives, anthoxanthins are also used in the production of dyes and pigments. Anthoxanthins can also be used to create yellow, orange, or red dyes for use in textiles, cosmetics, and other products.
| Physical sciences | Polyphenols | Chemistry |
2527306 | https://en.wikipedia.org/wiki/Patagona | Patagona | The giant hummingbirds are hummingbirds of the genus Patagona. The genus includes two species, the sedentary giant hummingbird and the migratory giant hummingbird, which are the largest and second largest species of hummingbird respectively.
Taxonomy
The giant hummingbird was described and illustrated in 1824 by the French ornithologist Louis Pierre Vieillot based on a specimen that Vieillot mistakenly believed had been collected in Brazil. The type locality was designated as Valparaíso in Chile by Carl Eduard Hellmayr in 1945. The giant hummingbird was the only species placed in the genus Patagona when introduced by George Robert Gray in 1840.
Molecular phylogenetic studies have shown that the giant hummingbird has no close relatives and is sister to the hummingbird subfamily Trochilinae, a large clade containing the tribes Lampornithini (mountain gems), Mellisugini (bees) and Trochilini (emeralds).
Two subspecies were previously recognised:
P. g. peruviana Boucard, 1893 – southwest Colombia to north Chile and northwest Argentina
P. g. gigas (Vieillot, 1824) – central, south Chile and west-central Argentina
These subspecies are thought to have emerged as a result of partial geographical separation of populations by volcanic activity in the Andes predating the Miocene; however, there remain areas of contact between the species, hence the lack of full speciation. The proposed phylogenetic system for hummingbirds suggested by McGuire et al. (2009) accommodated the possible elevation of these subspecies to species rank.
Description
Giant hummingbirds can be identified by their large size and characteristics such as the presence of an eye-ring, straight bill longer than the head, dull colouration, very long wings (approaching the tail tip when stowed), long and moderately forked tail, tarsi feathered to the toes and large, sturdy feet. There is no difference between the sexes. Juveniles have small corrugations on the lateral beak culmen.
Prior to the giant hummingbird being split into the Northern and Southern species, it was described as weighing and having a wingspan of approximately and length of . Although the Northern species is larger, both are approximately the same length as a European starling or a northern cardinal, though giant hummingbirds are considerably lighter because it has a slender build and long bill, making the body a smaller proportion of the total length. This weight is almost twice that of the heaviest hummingbird species outside of the genus Patagona and ten times that of the smallest, the bee hummingbird.
The giant hummingbird occasionally glides in flight, a behavior very rare among hummingbirds. Its elongated wings allow more efficient glides than do those of other hummingbirds. The giant hummingbird's voice is a distinctive loud, sharp and whistling "chip".
Behaviour
Hummingbirds are extremely agile and acrobatic flyers, regularly partaking in sustained hovering flight, often used not only to feed on the wing but to protect their territory and court mates. The giant hummingbird is typical in that it will brazenly defend its energy-rich flower territory from other species and other giant hummingbirds. These birds are typically seen alone, in pairs or small family groups.
Flight, anatomy and physiology
The giant hummingbird hovers at an average of 15 wing beats per second, a slow rate for a hummingbird. Its resting heart rate is 300 beats per minute, with a peak rate in flight of 1020 beats per minute. Energy requirements for hummingbirds do not scale evenly with size increases, meaning a larger bird such as giant hummingbird requires more energy per gram to hover than a smaller bird.
The giant hummingbird requires an estimated 4.3 calories of food energy per hour to sustain its flight. This requirement along with the low oxygen availability and thin air (generating little lift) at the high altitudes where the giant hummingbird usually lives suggest that it is close to the viable maximum size for a hummingbird.
Food and feeding
The giant hummingbird feeds mainly on nectar, visiting a range of flowers. The female giant hummingbird has been observed ingesting sources of calcium (sand, soil, slaked lime and wood ash) after the reproductive season to replenish the calcium used in egg production; the low calcium content of nectar necessitates these extra sources. Similarly, a nectar-based diet is low in protein and various dietary minerals, and this is countered by consuming insects.
It regularly feeds from the flowers of the genus Puya in Chile, with which it has a symbiotic relationship, trading pollination for food. As a large hovering bird, particularly at high altitudes, the giant hummingbird has extremely high metabolic requirements. It is known to feed from columnar cacti, including Oreocereus celsianus and Echinopsis atacamensis ssp. pasacana, and Salvia haenkei.
Considering the energy-rich nature of nectar as a food source, it attracts a large range of visitors apart from the hummingbird, which has coevolved with a plant to be the flower's most efficient pollinator. These other visitors, because they are not designed to access the well-hidden nectar, often damage the flowers (for example, piercing them at the base) and prevent further nectar production. Because of its high energy requirements, the giant hummingbird alters its foraging behaviour as a direct response to nectar robbing from other birds and animals, and this reduces the viability of the hummingbird in an area with many nectar robbers, as well as indirectly affecting the plants by reducing pollination. If alien species are introduced that become nectar thieves, it is reasonable to predict that their activities will significantly impact the local ecosystem. This could prove to be a future risk for the giant hummingbird populations because they sit close to the physical limit in their metabolic demands.
Breeding
There is little known of the giant hummingbird's breeding behaviour, but some generalisations can be inferred from other hummingbird species. Hummingbird males tend to have polygynous, occasionally promiscuous, behaviours, and no involvement after copulation. The female builds the nest and lays a clutch of two eggs during the summer. A giant hummingbird nest is small considering the size of the bird, typically made near water sources and perched on a branch of a tree or shrub parallel to the ground.
| Biology and health sciences | Apodiformes | Animals |
2527382 | https://en.wikipedia.org/wiki/Paleothermometer | Paleothermometer | A paleothermometer is a methodology that provides an estimate of the ambient temperature at the time of formation of a natural material. Most paleothermometers are based on empirically-calibrated proxy relationships, such as the tree ring or TEX86 methods. Isotope methods, such as the δ18O method or the clumped-isotope method, are able to provide, at least in theory, direct measurements of temperature.
Common paleothermometers
The isotopic ratio of 18O to 16O, usually in foram tests or ice cores. High values mean low temperatures. Confounded by ice volume - more ice means higher values.
Ocean water is mostly H216O, with small amounts of HD16O and H218O. In Standard Mean Ocean Water (SMOW) the ratio of D to H is and 18O/16O is . Fractionation occurs during changes between condensed and vapour phases: the vapour pressure of heavier isotopes is lower, so vapour contains relatively more of the lighter isotopes and when the vapour condenses the precipitation preferentially contains heavier isotopes. The difference from SMOW is expressed as
;
and a similar formula for δD. values for precipitation are always negative. The major influence on is the difference between ocean temperatures where the moisture evaporated and the place where the final precipitation occurred; since ocean temperatures are relatively stable the value mostly reflects the temperature where precipitation occurs. Taking into account that the precipitation forms above the inversion layer, we are left with a linear relation:
which is empirically calibrated from measurements of temperature and as a = for Greenland and for East Antarctica. The calibration was initially done on the basis of spatial variations in temperature and it was assumed that this corresponded to temporal variations (Jouzel and Merlivat, 1984). More recently, borehole thermometry has shown that for glacial-interglacial variations, a = (Cuffey et al., 1995), implying that glacial-interglacial temperature changes were twice as large as previously believed.
Mg/Ca and Sr/Ca
Magnesium (Mg) is incorporated into the calcite shells (tests) of planktic and benthic foraminifera as a trace element. Because the incorporation of Mg as an impurity in calcite is endothermic, more is incorporated into the growing crystal at higher temperatures. Therefore, a high Mg/Ca ratio implies a high temperature, although ecological factors may confound the signal. Mg has a long residence time in the ocean, and so it is possible to largely ignore the effect of changes in seawater Mg/Ca on the signal. Mg/Ca ratios can sometimes underestimate seawater temperatures by way of the dissolution of foraminifer shells, which lowers Mg/Ca values.
Strontium (Sr) incorporates in coral aragonite, and it is well established that the precise Sr/Ca ratio in the coral skeleton shows an inverse correlation with the seawater temperature during its biomineralization.
Alkenones
Distributions of organic molecules in marine sediments reflect temperature.
Leaf physiognomy
The characteristic leaf sizes, shapes and prevalence of features such as drip tips (‘leaf or foliar physiognomy’) differs between tropical rainforests (many species with large leaves with smooth edges and drip tips) and temperate deciduous forests (smaller leaf size classes common, toothed edges common), and is often continuously variable between sites along climatic gradients, such as from hot to cold climates, or high to low precipitation. This variation between sites along environmental gradients reflects adaptive compromises by the species present to balance the need to capture light energy, manage heat gain and loss, while maximising the efficiency of gas exchange, transpiration and photosynthesis. Quantitative analyses of modern vegetation leaf physiognomy and climate responses along environmental gradients have been largely univariate, but multivariate approaches integrate multiple leaf characters and climatic parameters. Temperature has been estimated (to varying degrees of fidelity) using leaf physiognomy for Late Cretaceous and Cenozoic leaf floras, principally using two main approaches:
Leaf margin analysis
A univariate approach that is based on the observation that the proportion of woody dicot species with smooth (i.e. non-toothed) leaf margins (0 ≤ Pmargin ≤ 1) in vegetation varies proportionately with mean annual temperature (MAT). Requires the fossil flora to be segregated into morphotypes (i.e. ‘species’), but does not require their identification. The original LMA regression equation was derived for East Asian forests, and is:
The error of the estimate for LMA is expressed as the binomial sampling error:
where c is the slope from the LMA regression equation, Pmargin as used in (), and r is the number of species scored for leaf margin type for the individual fossil leaf flora.
LMA calibrations have been derived for major world regions, including North America, Europe, South America, and Australia. Riparian and wetland environments have a slightly different regression equation, because they have proportionally fewer smooth-margined plants. It is
CLAMP (Climate leaf analysis multivariate program)
CLAMP is a multivariate approach largely based on a data set of primarily western hemisphere vegetation, subsequently added to with datasets from additional world regional vegetation. Canonical Correlation Analysis is used combining 31 leaf characters, but leaf margin type represented a significant component of the relationship between physiognomic states and temperature. Using CLAMP, MAT is estimated with small standard errors (e.g. CCA ± 0.7–1.0 °C). Additional temperature parameters can be estimated using CLAMP, such as the coldest month mean temperature (CMMT) and the warmest month mean temperature (WMMT) which provide estimates for winter and summer mean conditions respectively.
Nearest living relative analogy / coexistence analysis
Certain plants prefer certain temperatures; if their pollen is found one can work out the approximate temperature.
13C-18O bonds in carbonates
There is a slight thermodynamic tendency for heavy isotopes to form bonds with each other, in excess of what would be expected from a stochastic or random distribution of the same concentration of isotopes. The excess is greatest at low temperature (see Van 't Hoff equation), with the isotopic distribution becoming more randomized at higher temperature. Along with the closely related phenomenon of equilibrium isotope fractionation, this effect arises from differences in zero point energy among isotopologues. Carbonate minerals like calcite contain CO32− groups that can be converted to CO2 gas by reaction with concentrated phosphoric acid. The CO2 gas is analyzed with a mass spectrometer, to determine the abundances of isotopologues. The parameter Δ47 is the measured difference in concentration between isotopologues with a mass of 47 u (as compared to 44) in a sample and a hypothetical sample with the same bulk isotopic composition, but a stochastic distribution of heavy isotopes. Lab experiments, quantum mechanical calculations, and natural samples (with known crystallization temperatures) all indicate that Δ47 is correlated to the inverse square of temperature. Thus Δ47 measurements provide an estimation of the temperature at which a carbonate formed. 13C-18O paleothermometry does not require prior knowledge of the concentration of 18O in the water (which the δ18O method does). This allows the 13C-18O paleothermometer to be applied to some samples, including freshwater carbonates and very old rocks, with less ambiguity than other isotope-based methods. The method is presently limited by the very low concentration of isotopologues of mass 47 or higher in CO2 produced from natural carbonates, and by the scarcity of instruments with appropriate detector arrays and sensitivities. The study of these types of isotopic ordering reactions in nature is often called "clumped-isotope" geochemistry.
| Physical sciences | Paleoclimate | Earth science |
2528595 | https://en.wikipedia.org/wiki/Russet%20Burbank | Russet Burbank | Russet Burbank is a potato cultivar with dark brown skin and few eyes that is the most widely grown potato in North America. A russet type, its flesh is white, dry, and mealy, and it is good for baking, mashing, and french fries (chips). It is a common and popular potato.
Origin
This variety is a mutation (or sport) of the cultivar 'Burbank's Seedling' that was selected by the plant breeder Luther Burbank in 1873. The known lineage of Russet Burbank began in 1853 when Chauncey E. Goodrich imported the Rough Purple Chili from South America in an attempt to add diversity to American potato stocks which were susceptible to late blight. Goodrich bred Garnet Chili from Rough Purple Chili, and Albert Bresee bred Early Rose from Garnet Chili, from which Luther Burbank bred Burbank. This cross-over was formerly known as the Russell, but was eventually popularized as the Russet potato in the American stores. Russet Burbank has been widely, but incorrectly, reported to have been selected in 1914 by the Colorado potato grower Lou D. Sweet. A 2014 study confirmed that it was originally released in 1902 by L. L. May & Co and was first known as the Netted Gem.
To improve the disease resistance of Irish potatoes, Luther Burbank selected the potato that became known as "the Burbank." It was not patented because plants, such as potatoes, propagated from tubers were not granted patents in the United States.
Usage
Russet Burbank was not initially popular, accounting for only 4% of potatoes in the US in 1930. The introduction of irrigation in Idaho increased its popularity, as growers found it produced large potatoes easily marketed as baking potatoes. The invention of frozen french fries in the '40s and fast food restaurants in the '50s increased its popularity further. By the 2010s, Russet Burbank accounted for 70% of the ultra-processed potato market in North America, and over 40% of the potato growing area in the US.
Restaurants such as McDonald's favor russet potatoes for their size, which produce long pieces suitable for french fries. As of 2009, "McDonald's top tuber is the Russet Burbank." After decades of consumption in North America, consumers and processors consider it the standard potato against which others are judged.
Botanical features
The Russet Burbank plants are medium-sized with stems with a medium thickness that are prominently angled. The leaves of this variety are medium-sized with large terminal and primary leaflets. The plant's flowers are medium-sized with dark green buds that drop readily. The variety has large, long tubers that are cylindrical or slightly flat. There are numerous eyes on the potato that are evenly distributed, and the sprouts are brownish-purple.
Storage
The Russet Burbank variety stores very well for long periods. It can be stored at 7 °C (44.6 °F) for up to five months without the need to apply gasses that inhibit sprouting. One issue that can occur while in storage is an internal black spot, also known as IBS. Also, if the potatoes are harvested too early, there could be a skinning issue.
Disease resistance
Russet Burbank is highly resistant to black leg. It is moderately resistant to common scab and fusarium dry rot (Fusarium oxysporum and F. sambucinum). It is susceptible to fusarium dry rot (F. coeruleum), late blight (Phytophthora infestans), leaf roll, seed-piece decay, tuber net necrosis, verticillium wilt, PVX and PVY. Genetically modified potatoes can be resistant to the Colorado Potato Beetle.
| Biology and health sciences | Root vegetables | Plants |
2532789 | https://en.wikipedia.org/wiki/Standard%20gravity | Standard gravity | The standard acceleration of gravity or standard acceleration of free fall, often called simply standard gravity and denoted by or , is the nominal gravitational acceleration of an object in a vacuum near the surface of the Earth. It is a constant defined by standard as . This value was established by the third General Conference on Weights and Measures (1901, CR 70) and used to define the standard weight of an object as the product of its mass and this nominal acceleration. The acceleration of a body near the surface of the Earth is due to the combined effects of gravity and centrifugal acceleration from the rotation of the Earth (but the latter is small enough to be negligible for most purposes); the total (the apparent gravity) is about 0.5% greater at the poles than at the Equator.
Although the symbol is sometimes used for standard gravity, (without a suffix) can also mean the local acceleration due to local gravity and centrifugal acceleration, which varies depending on one's position on Earth (see Earth's gravity). The symbol should not be confused with , the gravitational constant, or g, the symbol for gram. The is also used as a unit for any form of acceleration, with the value defined as above.
The value of defined above is a nominal midrange value on Earth, originally based on the acceleration of a body in free fall at sea level at a geodetic latitude of 45°. Although the actual acceleration of free fall on Earth varies according to location, the above standard figure is always used for metrological purposes. In particular, since it is the ratio of the kilogram-force and the kilogram, its numeric value when expressed in coherent SI units is the ratio of the kilogram-force and the newton, two units of force.
History
Already in the early days of its existence, the International Committee for Weights and Measures (CIPM) proceeded to define a standard thermometric scale, using the boiling point of water. Since the boiling point varies with the atmospheric pressure, the CIPM needed to define a standard atmospheric pressure. The definition they chose was based on the weight of a column of mercury of 760 mm. But since that weight depends on the local gravity, they now also needed a standard gravity. The 1887 CIPM meeting decided as follows:
All that was needed to obtain a numerical value for standard gravity was now to measure the gravitational strength at the International Bureau. This task was given to Gilbert Étienne Defforges of the Geographic Service of the French Army. The value he found, based on measurements taken in March and April 1888, was 9.80991(5) m⋅s−2.
This result formed the basis for determining the value still used today for standard gravity. The third General Conference on Weights and Measures, held in 1901, adopted a resolution declaring as follows:
The numeric value adopted for was, in accordance with the 1887 CIPM declaration, obtained by dividing Defforges's result – 980.991 cm⋅s−2 in the cgs system then en vogue – by 1.0003322 while not taking more digits than are warranted considering the uncertainty in the result.
Conversions
| Physical sciences | Acceleration | Basics and measurement |
19469852 | https://en.wikipedia.org/wiki/Plane%20%28mathematics%29 | Plane (mathematics) | In mathematics, a plane is a two-dimensional space or flat surface that extends indefinitely.
A plane is the two-dimensional analogue of a point (zero dimensions), a line (one dimension) and three-dimensional space. When working exclusively in two-dimensional Euclidean space, the definite article is used, so the Euclidean plane refers to the whole space.
Several notions of a plane may be defined. The Euclidean plane follows Euclidean geometry, and in particular the parallel postulate. A projective plane may be constructed by adding "points at infinity" where two otherwise parallel lines would intersect, so that every pair of lines intersects in exactly one point. The elliptic plane may be further defined by adding a metric to the real projective plane. One may also conceive of a hyperbolic plane, which obeys hyperbolic geometry and has a negative curvature.
Abstractly, one may forget all structure except the topology, producing the topological plane, which is homeomorphic to an open disk. Viewing the plane as an affine space produces the affine plane, which lacks a notion of distance but preserves the notion of collinearity. Conversely, in adding more structure, one may view the plane as a 1-dimensional complex manifold, called the complex line.
Many fundamental tasks in mathematics, geometry, trigonometry, graph theory, and graphing are performed in a two-dimensional or planar space.
Euclidean plane
Embedding in three-dimensional space
Elliptic plane
Projective plane
Further generalizations
In addition to its familiar geometric structure, with isomorphisms that are isometries with respect to the usual inner product, the plane may be viewed at various other levels of abstraction. Each level of abstraction corresponds to a specific category.
At one extreme, all geometrical and metric concepts may be dropped to leave the topological plane, which may be thought of as an idealized homotopically trivial infinite rubber sheet, which retains a notion of proximity, but has no distances. The topological plane has a concept of a linear path, but no concept of a straight line. The topological plane, or its equivalent the open disc, is the basic topological neighborhood used to construct surfaces (or 2-manifolds) classified in low-dimensional topology. Isomorphisms of the topological plane are all continuous bijections. The topological plane is the natural context for the branch of graph theory that deals with planar graphs, and results such as the four color theorem.
The plane may also be viewed as an affine space, whose isomorphisms are combinations of translations and non-singular linear maps. From this viewpoint there are no distances, but collinearity and ratios of distances on any line are preserved.
Differential geometry views a plane as a 2-dimensional real manifold, a topological plane which is provided with a differential structure. Again in this case, there is no notion of distance, but there is now a concept of smoothness of maps, for example a differentiable or smooth path (depending on the type of differential structure applied). The isomorphisms in this case are bijections with the chosen degree of differentiability.
In the opposite direction of abstraction, we may apply a compatible field structure to the geometric plane, giving rise to the complex plane and the major area of complex analysis. The complex field has only two isomorphisms that leave the real line fixed, the identity and conjugation.
In the same way as in the real case, the plane may also be viewed as the simplest, one-dimensional (in terms of complex dimension, over the complex numbers) complex manifold, sometimes called the complex line. However, this viewpoint contrasts sharply with the case of the plane as a 2-dimensional real manifold. The isomorphisms are all conformal bijections of the complex plane, but the only possibilities are maps that correspond to the composition of a multiplication by a complex number and a translation.
In addition, the Euclidean geometry (which has zero curvature everywhere) is not the only geometry that the plane may have. The plane may be given a spherical geometry by using the stereographic projection. This can be thought of as placing a sphere tangent to the plane (just like a ball on the floor), removing the top point, and projecting the sphere onto the plane from this point. This is one of the projections that may be used in making a flat map of part of the Earth's surface. The resulting geometry has constant positive curvature.
Alternatively, the plane can also be given a metric which gives it constant negative curvature giving the hyperbolic plane. The latter possibility finds an application in the theory of special relativity in the simplified case where there are two spatial dimensions and one time dimension. (The hyperbolic plane is a timelike hypersurface in three-dimensional Minkowski space.)
Topological and differential geometric notions
The one-point compactification of the plane is homeomorphic to a sphere (see stereographic projection); the open disk is homeomorphic to a sphere with the "north pole" missing; adding that point completes the (compact) sphere. The result of this compactification is a manifold referred to as the Riemann sphere or the complex projective line. The projection from the Euclidean plane to a sphere without a point is a diffeomorphism and even a conformal map.
The plane itself is homeomorphic (and diffeomorphic) to an open disk. For the hyperbolic plane such diffeomorphism is conformal, but for the Euclidean plane it is not.
| Mathematics | Geometry | null |
19480112 | https://en.wikipedia.org/wiki/Arctic%20methane%20emissions | Arctic methane emissions | Arctic methane emissions contribute to a rise in methane concentrations in the atmosphere. Whilst the Arctic region is one of many natural sources of the greenhouse gas methane, there is nowadays also a human component to this due to the effects of climate change. In the Arctic, the main human-influenced sources of methane are thawing permafrost, Arctic sea ice melting, clathrate breakdown and Greenland ice sheet melting. This methane release results in a positive climate change feedback (meaning one that amplifies warming), as methane is a powerful greenhouse gas. When permafrost thaws due to global warming, large amounts of organic material can become available for methanogenesis and may therefore be released as methane.
Since around 2018, there has been consistent increases in global levels of methane in the atmosphere, with the 2020 increase of 15.06 parts per billion breaking the previous record increase of 14.05 ppb set in 1991, and 2021 setting an even larger increase of 18.34 ppb. However, there is currently no evidence connecting the Arctic to this recent acceleration. In fact, a 2021 study indicated that the methane contributions from the Arctic were generally overestimated, while the contributions of tropical regions were underestimated.
Nevertheless, the Arctic's role in global methane trends is considered very likely to increase in the future. There is evidence for increasing methane emissions since 2004 from a Siberian permafrost site into the atmosphere linked to warming.
Mitigation of CO2 emissions by 2050 (i.e. reaching net zero emissions) is probably not enough to stop the future disappearance of summer Arctic Ocean ice cover. Mitigation of methane emissions is also necessary and this has to be carried out over an even shorter period of time. Such mitigation activities need to be carried out in three main sectors: oil and gas, waste and agriculture. Using available measures this could amount to global reductions of ca.180 Mt/yr or about 45% of the current (2021) emissions by 2030.
Observed values and processes
NOAA annual records for methane concentrations in the atmosphere have been updated since 1984. They show substantial growth during the 1980s, a slowdown in annual growth during the 1990s, a plateau (including some years of declining atmospheric concentrations) in the early 2000s and another consistent increase beginning in 2007. Since around 2018, there has been consistent annual increases in global levels of methane, with the 2020 increase of 15.06 parts per billion breaking the previous record increase of 14.05 ppb set in 1991, and 2021 setting an even larger increase of 18.34 ppb.
Due to the relatively short lifetime of atmospheric methane (7-12 years compared to 100s of years for CO2) its global trends are more complex than those of carbon dioxide.
These trends alarm climate scientists, with some suggesting that they represent a climate change feedback increasing natural methane emissions well beyond their preindustrial levels. However, there is currently no evidence connecting the Arctic to this recent acceleration. In fact, a 2021 study indicated that the role of the Arctic was typically overestimated in global methane accounting, while the role of tropical regions was consistently underestimated. The study suggested that tropical wetland methane emissions were the culprit behind the recent growth trend, and this hypothesis was reinforced by a 2022 paper connecting tropical terrestrial emissions to 80% of the global atmospheric methane trends between 2010 and 2019.
Nevertheless, the Arctic's role in global methane trends is considered very likely to increase in the future. There is evidence for increasing methane emissions since 2004 from a Siberian permafrost site into the atmosphere linked to warming.
Radiocarbon dating of trace methane in lake bubbles and soil organic carbon concluded that 0.2 to 2.5 Pg of permafrost carbon has been released as methane and carbon dioxide over the last 60 years. The 2020 heat wave may have released significant methane from carbonate deposits in Siberian permafrost.
Methane emissions by the permafrost carbon feedback—amplification of surface warming due to enhanced radiative forcing by carbon release from permafrost—could contribute an estimated 205 Gt of carbon emissions, leading up to 0.5 °C (0.9 °F) of additional warming by the end of the 21st century. However, recent research based on the carbon isotopic composition of atmospheric methane trapped in bubbles in Antarctic ice suggests that methane emissions from permafrost and methane hydrates were minor during the last deglaciation, suggesting that future permafrost methane emissions may be lower than previously estimated.
Comparison of Arctic and Antarctic atmosphere measurements
Atmospheric methane concentrations are 8–10% higher in the Arctic than in the Antarctic atmosphere. During cold glacier epochs, this gradient decreases to insignificant levels. Land ecosystems are thought to be the main sources of this asymmetry, although it has been suggested in 2007 that "the role of the Arctic Ocean is significantly underestimated." Soil temperature and moisture levels are important variables in soil methane fluxes in tundra environments.
Sources of methane in the Arctic
Large quantities of methane are stored in the Arctic in natural gas deposits, permafrost, and as undersea clathrates. Permafrost and clathrates degrade on warming, thus large releases of methane from these sources may arise as a result of global warming. Other sources of methane include submarine taliks, river transport, ice complex retreat, submarine permafrost and decaying gas hydrate deposits. Permafrost contains almost twice as much carbon as the atmosphere, with ~20 Gt of permafrost-associated methane trapped in methane clathrates. Permafrost thaw results in the formation of thermokarst lakes in ice-rich yedoma deposits. Methane frozen in permafrost is slowly released as permafrost thaws.
Thawing permafrost
Arctic sea ice melting
Clathrate breakdown
Greenland ice sheet melting
A 2014 study found evidence for methane cycling below the ice sheet of the Russell Glacier, based on subglacial drainage samples which were dominated by Pseudomonadota bacteria. During the study, the most widespread surface melt on record for the past 120 years was observed in Greenland; on 12 July 2012, unfrozen water was present on almost the entire ice sheet surface (98.6%). The findings indicate that methanotrophs could serve as a biological methane sink in the subglacial ecosystem, and the region was, at least during the sample time, a source of atmospheric methane. Scaled dissolved methane flux during the four months of the summer melt season for the Russell Glacier catchment area (1200 km2) was estimated at 990 tonnes CH4. Because this catchment area is representative of similar Greenland outlet glaciers, the researchers concluded that the Greenland Ice Sheet may represent a significant global methane source.
A study in 2016 concluded that methane clathrates may exist below Greenland's and Antarctica's ice sheets, based on past evidence.
Reducing methane emissions
More than half of global methane emissions originate from human activities across three main sectors: fossil fuels (35% of human-caused emissions), waste (20%), and agriculture (40%). Within the fossil fuel sector, oil and gas extraction, processing, and distribution contribute 23%, while coal mining accounts for 12% of these emissions. In the waste sector, landfills and wastewater comprise about 20% of global anthropogenic emissions. In agriculture, livestock emissions from manure and enteric fermentation make up roughly 32%, and rice cultivation contributes 8% of global anthropogenic emissions. Mitigation using available measures could reduce these methane emissions by about 180 Mt/yr or about 45% by 2030.
Mitigation of CO2 emissions by 2050 (i.e. reaching net zero emissions) is probably not enough to stop the future disappearance of summer Arctic Ocean ice cover. Mitigation of methane emissions is also necessary and this has to be carried out over an even shorter period of time.
Flaring methane from oil and gas operations
ARPA-E has funded a research project from 2021-2023 to develop a "smart micro-flare fleet" to burn off methane emissions at remote locations.
A 2012 review article stated that most existing technologies "operate on confined gas streams of 0.1% methane", and were most suitable for areas where methane is emitted in pockets.
If Arctic oil and gas operations use Best Available Technology (BAT) and Best Environmental Practices (BEP) in petroleum gas flaring, this can result in significant methane emissions reductions, according to the Arctic Council.
| Physical sciences | Climate change | Earth science |
19480890 | https://en.wikipedia.org/wiki/Law%20of%20cosines | Law of cosines | In trigonometry, the law of cosines (also known as the cosine formula or cosine rule) relates the lengths of the sides of a triangle to the cosine of one of its angles. For a triangle with sides and opposite respective angles and (see Fig. 1), the law of cosines states:
The law of cosines generalizes the Pythagorean theorem, which holds only for right triangles: if is a right angle then and the law of cosines reduces to
The law of cosines is useful for solving a triangle when all three sides or two sides and their included angle are given.
Use in solving triangles
The theorem is used in solution of triangles, i.e., to find (see Figure 3):
the third side of a triangle if two sides and the angle between them is known:
the angles of a triangle if the three sides are known:
the third side of a triangle if two sides and an angle opposite to one of them is known (this side can also be found by two applications of the law of sines):
These formulas produce high round-off errors in floating point calculations if the triangle is very acute, i.e., if is small relative to and or is small compared to 1. It is even possible to obtain a result slightly greater than one for the cosine of an angle.
The third formula shown is the result of solving for a in the quadratic equation . This equation can have 2, 1, or 0 positive solutions corresponding to the number of possible triangles given the data. It will have two positive solutions if , only one positive solution if , and no solution if . These different cases are also explained by the side-side-angle congruence ambiguity.
History
Book II of Euclid's Elements, compiled c. 300 BC from material up to a century or two older, contains a geometric theorem corresponding to the law of cosines but expressed in the contemporary language of rectangle areas; Hellenistic trigonometry developed later, and sine and cosine per se first appeared centuries afterward in India.
The cases of obtuse triangles and acute triangles (corresponding to the two cases of negative or positive cosine) are treated separately, in Propositions II.12 and II.13:
Proposition 13 contains an analogous statement for acute triangles. In his (now-lost and only preserved through fragmentary quotations) commentary, Heron of Alexandria provided proofs of the converses of both II.12 and II.13.
Using notation as in Fig. 2, Euclid's statement of proposition II.12 can be represented more concisely (though anachronistically) by the formula
To transform this into the familiar expression for the law of cosines, substitute and
Proposition II.13 was not used in Euclid's time for the solution of triangles, but later it was used that way in the course of solving astronomical problems by al-Bīrūnī (11th century) and Johannes de Muris (14th century). Something equivalent to the spherical law of cosines was used (but not stated in general) by al-Khwārizmī (9th century), al-Battānī (9th century), and Nīlakaṇṭha (15th century).
The 13th century Persian mathematician Naṣīr al-Dīn al-Ṭūsī, in his (Book on the Complete Quadrilateral, c. 1250), gave a method for finding the third side of a general scalene triangle given two sides and the included angle by dropping a perpendicular from the vertex of one of the unknown angles to the opposite base, reducing the problem to solving a right-angled triangle by the Pythagorean theorem. If written out using modern mathematical notation, the resulting relation can be algebraically manipulated into the modern law of cosines.
About two centuries later, another Persian mathematician, Jamshīd al-Kāshī, who computed the most accurate trigonometric tables of his era, wrote about various methods of solving triangles in his (Key of Arithmetic, 1427), and repeated essentially al-Ṭūsī's method, including more explicit details, as follows:
Using modern algebraic notation and conventions this might be written
when is acute or
when is obtuse. (When is obtuse, the modern convention is that is negative and is positive; historically sines and cosines were considered to be line segments with non-negative lengths.) By squaring both sides, expanding the squared binomial, and then applying the Pythagorean trigonometric identity we obtain the familiar law of cosines:
In France, the law of cosines is sometimes referred to as the théorème d'Al-Kashi.
The theorem was first written using algebraic notation by François Viète in the 16th century. At the beginning of the 19th century, modern algebraic notation allowed the law of cosines to be written in its current symbolic form.
Proofs
Using the Pythagorean theorem
Case of an obtuse angle
Euclid proved this theorem by applying the Pythagorean theorem to each of the two right triangles in Fig. 2 ( and ). Using to denote the line segment and for the height , triangle gives us
and triangle gives
Expanding the first equation gives
Substituting the second equation into this, the following can be obtained:
This is Euclid's Proposition 12 from Book 2 of the Elements. To transform it into the modern form of the law of cosines, note that
Case of an acute angle
Euclid's proof of his Proposition 13 proceeds along the same lines as his proof of Proposition 12: he applies the Pythagorean theorem to both right triangles formed by dropping the perpendicular onto one of the sides enclosing the angle and uses the square of a difference to simplify.
Another proof in the acute case
Using more trigonometry, the law of cosines can be deduced by using the Pythagorean theorem only once. In fact, by using the right triangle on the left hand side of Fig. 6 it can be shown that:
using the trigonometric identity
This proof needs a slight modification if . In this case, the right triangle to which the Pythagorean theorem is applied moves outside the triangle . The only effect this has on the calculation is that the quantity is replaced by As this quantity enters the calculation only through its square, the rest of the proof is unaffected. However, this problem only occurs when is obtuse, and may be avoided by reflecting the triangle about the bisector of .
Referring to Fig. 6 it is worth noting that if the angle opposite side is then:
This is useful for direct calculation of a second angle when two sides and an included angle are given.
From three altitudes
The altitude through vertex is a segment perpendicular to side . The distance from the foot of the altitude to vertex plus the distance from the foot of the altitude to vertex is equal to the length of side (see Fig. 5). Each of these distances can be written as one of the other sides multiplied by the cosine of the adjacent angle,
(This is still true if or is obtuse, in which case the perpendicular falls outside the triangle.) Multiplying both sides by yields
The same steps work just as well when treating either of the other sides as the base of the triangle:
Taking the equation for and subtracting the equations for and
This proof is independent of the Pythagorean theorem, insofar as it is based only on the right-triangle definition of cosine and obtains squared side lengths algebraically. Other proofs typically invoke the Pythagorean theorem explicitly, and are more geometric, treating as a label for the length of a certain line segment.
Unlike many proofs, this one handles the cases of obtuse and acute angles in a unified fashion.
Cartesian coordinates
Consider a triangle with sides of length , , , where is the measurement of the angle opposite the side of length . This triangle can be placed on the Cartesian coordinate system with side aligned along the x axis and angle placed at the origin, by plotting the components of the 3 points of the triangle as shown in Fig. 4:
By the distance formula,
Squaring both sides and simplifying
An advantage of this proof is that it does not require the consideration of separate cases depending on whether the angle is acute, right, or obtuse. However, the cases treated separately in Elements II.12–13 and later by al-Ṭūsī, al-Kāshī, and others could themselves be combined by using concepts of signed lengths and areas and a concept of signed cosine, without needing a full Cartesian coordinate system.
Using Ptolemy's theorem
Referring to the diagram, triangle ABC with sides = , = and = is drawn inside its circumcircle as shown. Triangle is constructed congruent to triangle with = and = . Perpendiculars from and meet base at and respectively. Then:
Now the law of cosines is rendered by a straightforward application of Ptolemy's theorem to cyclic quadrilateral :
Plainly if angle is right, then is a rectangle and application of Ptolemy's theorem yields the Pythagorean theorem:
By comparing areas
One can also prove the law of cosines by calculating areas. The change of sign as the angle becomes obtuse makes a case distinction necessary.
Recall that
, , and are the areas of the squares with sides , , and , respectively;
if is acute, then is the area of the parallelogram with sides and forming an angle of ;
if is obtuse, and so is negative, then is the area of the parallelogram with sides a and b forming an angle of .
Acute case. Figure 7a shows a heptagon cut into smaller pieces (in two different ways) to yield a proof of the law of cosines. The various pieces are
in pink, the areas , on the left and the areas and on the right;
in blue, the triangle , on the left and on the right;
in grey, auxiliary triangles, all congruent to , an equal number (namely 2) both on the left and on the right.
The equality of areas on the left and on the right gives
Obtuse case. Figure 7b cuts a hexagon in two different ways into smaller pieces, yielding a proof of the law of cosines in the case that the angle is obtuse. We have
in pink, the areas , , and on the left and on the right;
in blue, the triangle twice, on the left, as well as on the right.
The equality of areas on the left and on the right gives
The rigorous proof will have to include proofs that various shapes are congruent and therefore have equal area. This will use the theory of congruent triangles.
Using circle geometry
Using the geometry of the circle, it is possible to give a more geometric proof than using the Pythagorean theorem alone. Algebraic manipulations (in particular the binomial theorem) are avoided.
Case of acute angle , where . Drop the perpendicular from onto = , creating a line segment of length . Duplicate the right triangle to form the isosceles triangle . Construct the circle with center and radius , and its tangent through . The tangent forms a right angle with the radius (Euclid's Elements: Book 3, Proposition 18; or see here), so the yellow triangle in Figure 8 is right. Apply the Pythagorean theorem to obtain
Then use the tangent secant theorem (Euclid's Elements: Book 3, Proposition 36), which says that the square on the tangent through a point outside the circle is equal to the product of the two lines segments (from ) created by any secant of the circle through . In the present case: , or
Substituting into the previous equation gives the law of cosines:
Note that is the power of the point with respect to the circle. The use of the Pythagorean theorem and the tangent secant theorem can be replaced by a single application of the power of a point theorem.
Case of acute angle , where . Drop the perpendicular from onto = , creating a line segment of length . Duplicate the right triangle to form the isosceles triangle . Construct the circle with center and radius , and a chord through perpendicular to half of which is Apply the Pythagorean theorem to obtain
Now use the chord theorem (Euclid's Elements: Book 3, Proposition 35), which says that if two chords intersect, the product of the two line segments obtained on one chord is equal to the product of the two line segments obtained on the other chord. In the present case: or
Substituting into the previous equation gives the law of cosines:
Note that the power of the point with respect to the circle has the negative value .
Case of obtuse angle . This proof uses the power of a point theorem directly, without the auxiliary triangles obtained by constructing a tangent or a chord. Construct a circle with center and radius (see Figure 9), which intersects the secant through and in and . The power of the point with respect to the circle is equal to both and . Therefore,
which is the law of cosines.
Using algebraic measures for line segments (allowing negative numbers as lengths of segments) the case of obtuse angle () and acute angle () can be treated simultaneously.
Using the law of sines
The law of cosines can be proven algebraically from the law of sines and a few standard trigonometric identities. To start, three angles of a triangle sum to a straight angle ( radians). Thus by the angle sum identities for sine and cosine,
Squaring the first of these identities, then substituting from the second, and finally replacing the Pythagorean trigonometric identity, we have:
The law of sines holds that
so to prove the law of cosines, we multiply both sides of our previous identity by
This concludes the proof.
Using vectors
Denote
Therefore,
Taking the dot product of each side with itself:
Using the identity
leads to
The result follows.
Isosceles case
When , i.e., when the triangle is isosceles with the two sides incident to the angle equal, the law of cosines simplifies significantly. Namely, because , the law of cosines becomes
or
Analogue for tetrahedra
Given an arbitrary tetrahedron whose four faces have areas , , , and , with dihedral angle between faces and , etc., a higher-dimensional analogue of the law of cosines is:
Version suited to small angles
When the angle, , is small and the adjacent sides, and , are of similar length, the right hand side of the standard form of the law of cosines is subject to catastrophic cancellation in numerical approximations. In situations where this is an important concern, a mathematically equivalent version of the law of cosines, similar to the haversine formula, can prove useful:
In the limit of an infinitesimal angle, the law of cosines degenerates into the circular arc length formula, .
In non-Euclidean geometry
As in Euclidean geometry, one can use the law of cosines to determine the angles , , from the knowledge of the sides , , . In contrast to Euclidean geometry, the reverse is also possible in both non-Euclidean models: the angles , , determine the sides , , .
A triangle is defined by three points , , and on the unit sphere, and the arcs of great circles connecting those points. If these great circles make angles , , and with opposite sides , , then the spherical law of cosines asserts that all of the following relationships hold:
In hyperbolic geometry, a pair of equations are collectively known as the hyperbolic law of cosines. The first is
where and are the hyperbolic sine and cosine, and the second is
The length of the sides can be computed by:
Polyhedra
The Law of Cosines can be generalized to all polyhedra by considering any polyhedron with vector sides and invoking the Divergence Theorem.
| Mathematics | Trigonometry | null |
10428984 | https://en.wikipedia.org/wiki/Tourniquet | Tourniquet | A tourniquet is a device that is used to apply pressure to a limb or extremity in order to create ischemia or stopping the flow of blood. It may be used in emergencies, in surgery, or in post-operative rehabilitation.
A simple tourniquet can be made from a stick and a rope, but the use of makeshift tourniquets has been reduced over time due to their ineffectiveness compared to a commercial and professional tourniquet. This may stem the flow of blood, but side effects such as soft tissue damage and nerve damage may occur.
History
During Alexander the Great’s military campaigns in the fourth century BC, tourniquets were used to stanch the bleeding of wounded soldiers. Romans used them to control bleeding, especially during amputations. These tourniquets were narrow straps made of bronze, using only leather for comfort.
In 1718, French surgeon Jean Louis Petit developed a screw device for occluding blood flow in surgical sites. Before this invention, the tourniquet was a simple garrot, tightened by twisting a rod (thus its name tourniquet, from tourner = to turn).
In 1785, Sir Gilbert Blane advocated that, in battle, each Royal Navy sailor should carry a tourniquet:
It frequently happens that men bleed to death before assistance can be procured, or lose so much blood as not to be able to go through an operation. In order to prevent this, it has been proposed, and on some occasions practised, to make each man carry about him a garter, or piece of rope yarn, in order to bind up a limb in case of profuse bleeding. If it be objected, that this, from its solemnity may be apt to intimidate common men, officers at least should make use of some precaution, especially as many of them, and those of the highest rank, are stationed on the quarter deck, which is one of the most exposed situations, and far removed from the cockpit, where the surgeon and his assistants are placed. This was the cause of the death of my friend Captain Bayne, of the Alfred, who having had his knee so shattered with round shot that it was necessary to amputate the limb, expired under the operation, in consequence of the weakness induced by loss of blood in carrying him so far. As the Admiral on these occasions allowed me the honour of being at his side, I carried in my pocket several tourniquets of a simple construction, in case that accidents to any person on the quarter deck should have required their use.
In 1864, Joseph Lister created a bloodless surgical field using a tourniquet device. In 1873, Friedrich von Esmarch introduced a rubber bandage that would both control bleeding and exsanguinate. This device is known as Esmarch's bandage. In 1881, Richard von Volkmann noted paralysis can occur from the use of the Esmarch tourniquet, if wrapped too tightly. Many cases of serious and permanent limb paralysis were reported from the use of non-pneumatic Esmarch tourniquets.
After observing considerable number of pressure paralysis with non-pneumatic, elastic, tourniquets, Harvey Cushing created a pneumatic tourniquet, in 1904. Pneumatic tourniquets were superior over Esmarch’s tourniquet in two ways: (1) faster application and removal; and (2) decrease the risk of nerve palsy.
In 1908, August Bier used two pneumatic tourniquets with intravenous local anesthesia to anesthetize the limb without general anesthetics.
In the early 1980s, microprocessor-based pneumatic tourniquet systems were invented by James McEwen. These modern electronic pneumatic tourniquet systems generally regulate the pressure in the tourniquet cuff within 1% of the target pressure and allows real-time monitoring of the inflation time. Modern pneumatic tourniquet systems include audiovisual alarms to alarm the user if hazardously high or low cuff pressures are present, automatic self-test and calibration, and backup power source.
In the 2000s, the silicon ring tourniquet, or elastic ring tourniquet, was developed by Noam Gavriely, a professor of medicine and former emergency physician. The tourniquet consists of an elastic ring made of silicone, stockinet, and pull straps made from ribbon that are used to roll the device onto the limb. The silicone ring tourniquet exsanguinates the blood from the limb while the device is being rolled on, and then occludes the limb once the desired occlusion location is reached. Unlike the historical mechanical tourniquets, the device reduces the risk of nerve paralysis. The surgical tourniquet version of the device is completely sterile, and provides improved surgical accessibility due to its narrow profile that results in a larger surgical field. It has been found to be a safe alternative method for most orthopedic limb procedures, but it does not completely replace the use of contemporary tourniquet devices. More recently the silicone ring tourniquet has been used in the fields of emergency medicine and vascular procedures. However, in 2015 Feldman et. al. reported two cases of pulmonary embolism after silicon ring exsanguination tourniquet application in patients with traumatic injuries. In one case of exsanguination tourniquet induced bilateral pulmonary emboli, after rapid intervention a 65-year-old woman was discharged in good condition 7 days after surgery. In a second case with multiple pulmonary emboli, despite extensive efforts of intervention a 53-year-old man’s condition quickly deteriorated after surgery, and was declared brain dead 2 days after. While Feldman et. al. discuss the potential risk of DVT for various types of tourniquets and exsanguination methods, the authors recommend extreme caution and suggest avoiding the use of an exsanguination tourniquet in patients with risk factors for DVT, including patients with traumatic injury of the extremities.
Most modern pneumatic tourniquet systems include the ability to measure the patient’s limb occlusion pressure (LOP) and recommend a tourniquet pressure based on the measured LOP to set safer and lower tourniquet pressures. Limb occlusion pressure is defined as "the minimum pressure required, at a specific time by a specific tourniquet cuff applied to a specific patient’s limb at a specific location, to stop the flow of arterial blood into the limb distal to the cuff.”
After World War II, the US military reduced use of the tourniquet because the time between application and reaching medical attention was so long that the damage from stopped circulation was worse than that from blood loss.
Since the beginning of the 21st century, US authorities have resuscitated its use in both military and non-military situations because treatment delays have been dramatically reduced. The Virginia State Police and police departments in Dallas, Philadelphia and other major cities provide tourniquets and other advanced bandages. In Afghanistan and Iraq, only 2 percent of soldiers with severe bleeding died compared with 7 percent in the Vietnam War, in part because of the combination of tourniquets and rapid access to doctors. Between 2005 and 2011, tourniquets saved 2,000 American lives from the wars in Iraq and Afghanistan. In civilian use, emerging practices include transporting tourniquetted patients even before emergency responders arrive and including tourniquets with defibrillators for emergency use.
There are currently no standards for testing tourniquets although there have been several proposed devices to ensure that the appropriate pressures could be generated including many commercial systems and an open source system that can be largely 3D printed. This would allow distributed manufacturing of tourniquets.
Risks
Risks and contraindications related to the use of a surgical tourniquet include: nerve injuries, skin injuries, compartment syndrome, deep venous thrombosis, and pain. Risk of injury can be minimized by minimizing tourniquet pressure and pressure gradients. Tourniquet pressure and pressure gradients can be minimized by using a tourniquet pressure based on the patient’s limb occlusion pressure, and by using a wider, contoured pneumatic tourniquet cuff.
In some elective surgical procedures such as total knee arthroplasty, some research suggests tourniquet use may be associated with an increased risk of adverse events, pain, and a longer hospital stay, despite tourniquet use allowing shorter times in the operating room. However, such evidence (meta-analyses and reviews) often omit the analysis of key tourniquet parameters and their correlation to outcomes leading to limited, inconclusive, and conflicting results.
A study by Pavao et al compared no tourniquet use to optimized tourniquet use in total knee arthroplasty and found no significant differences in surgical timing, blood loss, thigh and knee pain, edema, range of motion, functional scores, and complications, thus allowing surgery to occur with the benefits of a clean and dry surgical field from an optimized tourniquet without increase procedure-related comorbidities. Therefore, tourniquet use optimized to mitigate tourniquet related-risks while maintaining the benefits of a clear bloodless field and faster operating times may be achieved by minimizing tourniquet pressure and inflated tourniquet times.
Types
There are three types of tourniquets: surgical tourniquets, emergency tourniquets, and rehabilitation tourniquets.
Surgical tourniquets
Surgical tourniquets prevent blood flow to a limb and enable surgeons to work in a bloodless operative field. This allows surgical procedures to be performed with improved precision, safety and speed. Surgical tourniquets can be divided into two groups: pneumatic tourniquets and non-pneumatic tourniquets.
Surgical pneumatic tourniquets
Surgical pneumatic tourniquets are routinely and safely used orthopedic and plastic surgery, as well as in intravenous regional anesthesia (Bier block anesthesia) where they serve the additional function of preventing the central spread of local anesthetics in the limb. Modern pneumatic tourniquet systems consist of a pneumatic tourniquet instrument, tourniquet cuffs, pneumatic tubing, and limb protection sleeves.
Surgical pneumatic tourniquet instrument
Modern pneumatic tourniquet instruments are microcomputer-based with the following features:
Accurate pressure regulator to maintain cuff pressure within 1% of the target pressure,
Automatic timer to provide precise record of inflation time,
Audiovisual alarms to warn the operator if potential hazards are detected,
Automatic self test and self-calibration to ensure system hardware and software integrity, and
Backup power source to allow continued operation if unanticipated power outage occurs
Many studies published in the medical literature have shown that higher tourniquet pressures and pressure gradients are associated with higher risks of tourniquet-related injuries. Advances in tourniquet technology have reduced the risk of nerve-related injury by optimizing and personalizing tourniquet pressure based on the patient’s Limb Occlusion Pressure (LOP), rather than setting standard tourniquet pressures, which are generally higher and more hazardous. LOP is defined as “the minimum pressure required, at a specific time by a specific tourniquet cuff applied to a specific patient’s limb at a specific location, to stop the flow of arterial blood into the limb distal to the cuff.” LOP accounts for variables such as cuff design (bladder width), cuff application (snugness), patient limb characteristics (shape, size, tissues), and patient’s systolic blood pressure. After LOP is measured, personalized tourniquet pressure is set to LOP plus a safety margin to account for any increase in limb occlusion pressure normally expected during the surgery. The use of personalized pressures and wide contour tourniquet cuffs have been found to reduce average tourniquet pressure by 33%-42% from typical pressures. Setting the tourniquet pressure on the basis of LOP minimizes the pressure and related pressure gradients applied by a cuff to an underlying limb, which helps to minimize the risk of tourniquet-related injuries.
LOP may be measured manually by Doppler ultrasound. However, the method is time consuming and its accuracy is highly dependent on the skill and experience of the operator. LOP may also be measured automatically using a photoplethysmography distal sensor applied to the patient’s finger or toe of the operative limb to detect volumetric changes in blood in peripheral circulation as cuff pressure is gradually increased. Finally, most recently, LOP may be measured using a dual-purpose tourniquet cuff to monitor arterial pulsations in the underlying limb as the cuff pressure is gradually increased.
Pneumatic tourniquet instruments and cuffs are available in a single-line (single-port) or dual-line (dual-port) setup. Single-port configuration uses the same pneumatic line that connects the instrument to the cuff for both pressure regulation and pressure monitoring. Dual-port configuration uses one pneumatic line to regulate pressure and one pneumatic line to monitor pressure. The dual-port configuration may facilitate faster cuff pressure regulation and the detection of occlusions in the hoses.
Surgical pneumatic tourniquet cuff
Compressed gas is introduced into a bladder within a pneumatic tourniquet cuff by the pneumatic tourniquet instrument through a pneumatic tubing. The inflated cuff exerts pressure on the circumference of the patient’s limb to occlude blood flow.
Compression by the inflated cuff can result in tissue injury. A good tourniquet cuff fit ensures even pressure distribution across the underlying soft tissues, whereas a poor tourniquet cuff fit can result in areas of higher pressure which can lead to soft tissue ischemia. Therefore, in order to safely and effectively occlude blood flow distal to the applied tourniquet cuff, proper selection and application of the tourniquet cuff should be followed.
The following should be considered when selecting a tourniquet cuff:
Cuff location,
Limb shape which determines the cuff shape (e.g. cylindrical or contour shaped),
Limb circumference which determines the cuff length,
Cuff width,
Single versus dual bladder design (e.g. whether an IVRA cuff is needed), and
Use sterile cuff when it will be very close to the sterile field
Surgical limb protection sleeve
It is recommended to protect the limb beneath the cuff by applying a low-lint, soft padding around the limb, prior to cuff application, according to the cuff manufacturer’s instructions for use. Matching limb protection sleeves matched to the cuff width and patient’s limb circumference has been shown to produce significantly fewer, less severe wrinkles and pinches in the skin surface than other padding types tested.
Surgical non-pneumatic tourniquet
In silicone ring tourniquets, or elastic ring tourniquets, the tourniquet comes in a variety of sizes. To determine the correct tourniquet size, the patient's limb circumference at the desired occlusion location should be measured, as well as their blood pressure to determine the best model. Once the correct model is selected, typically two sterile medical personnel will be needed to apply the device. Unlike with a pneumatic tourniquet, the silicone ring tourniquet should be applied after the drapes have been placed on the patient. This is due to the device being completely sterile. The majority of the devices require a two-man operation (with the exception of the extra large model):
One person is responsible for holding the patient's limb. The other will place the device on the limb (extra large models may require two people).
Application:
The elastic ring tourniquet is placed on the patient's limb. If placed on a hand or foot, all fingers or toes should be enclosed within the tourniquet.
The handles of the tourniquet should be positioned medial-lateral on the upper extremity or posterior-anterior on the lower extremity.
The person applying the device should start rolling the device while the individual responsible for the limb should hold the limb straight and maintain axial traction.
Once the desired occlusion location is reached, the straps can be cut off or tied just below the ring.
A window can be cut or the section of stockinet can be completely removed.
Once the surgery is completed the device is cut off with a supplied cutting card.
The elastic ring tourniquet follows similar recommendations noted for pneumatic tourniquet use:
It should not be used on a patient's limb for more than 120 minutes, as the interruption of blood flow may cause cell damage and necrosis.
The tourniquet should not be placed on the ulnar nerve or the peroneal nerve.
The silicone ring device cannot be used on patients with blood problems such as DVT, edema, etc.
A patient suffering from skin lesions or a malignancy should use this type of tourniquet.
Emergency tourniquets
Emergency tourniquets differ from surgical tourniquets as are they are used in military combat care, emergency medicine, and accident situations where electrical power is not available, and may need to be applied by an assisting person or self-applied by the injured person. Emergency tourniquets are assessed for their effectiveness of hemorrhage control, pulse stoppage distal to the tourniquet, time to stop bleeding, total blood loss, and applied pressure. However, their design and safe use should be considered as it relates to nerve injury, reperfusion injury, soft tissue injury, and pain.
Early implementation of non-pneumatic tourniquet use in the nineteenth century for non-amputation surgical procedures often resulted in reports of permanent and temporary limb paralysis, nerve injuries, and other soft-tissue injuries. As a result, pneumatic tourniquets were developed for surgery, where the applied pressure and pressure gradients can be controlled, minimized, and controlled, and thereby minimize the risk of tourniquet related injuries.
Pneumatic emergency tourniquet
Emergency military tourniquet
The Emergency & Military Tourniquet (EMT) is an example of a pneumatic tourniquet developed for safe use in pre-hospital or military settings. In a study that evaluated 5 emergency tourniquet systems for use in the Canadian Forces, the EMT was one of the most effective tourniquets and caused the least pain. In another study comparing the effectiveness of 3 emergency tourniquet systems, while all devices were effective in both hemorrhage control and stopping blood flow, the EMT also performed the best for shortest time to stop blood flow, lowest total blood loss, and required the least amount of pressure to stop blood flow.
Non-pneumatic emergency tourniquet
Silicone ring auto-transfusion tourniquet
The silicone ring auto-transfusion tourniquet (SRT/ATT/EED), or surgical auto-transfusion tourniquet (HemaClear), is a simple to use, self-contained, mechanical tourniquet that consists of a silicone ring, stockinet, and pull straps that results in the limb being exsanguinated and occluded within seconds of application. The tourniquet can be used for limb procedures in the operating room, or in emergency medicine as a means to stabilize a patient until further treatment can be applied.
Combat application tourniquet
The combat application tourniquet (CAT) was developed by Ted Westmoreland. It is used by the U.S. and coalition militaries to provide soldiers a small and effective tourniquet in field combat situations. It is also used in the UK by NHS ambulance services, along with some UK fire and rescue services. The unit utilizes a windlass with a locking mechanism and can be self-applied. The CAT has been adopted by military and emergency personnel around the world.
An open hardware-based 3D printing project called the Glia Tourniquet (windlass type) enables emergency tourniquets to use distributed manufacturing to make them for $7 in materials. Concerns over quality control of distributed manufactured tourniquets was partially addressed with an open source testing apparatus. The tourniquet tester costs less than $100 and once calibrated with a blood pressure monitor, the built-in LCD displays the measuring range of the tester (0 to 200 N), which can be used to test the validation of all tourniquets.
Rehabilitation tourniquets
Personalized blood flow restriction
Recently, pneumatic tourniquets have been successfully used for a technique called Personalized Blood Flow Restriction Training (PBFRT) to accelerate the rehabilitation of orthopedic patients, injured professional athletes, and wounded soldiers.
Typically, to increase muscle size and strength, a person needs to lift loads at or above 65% of their one repetition maximum. However, injured patients are often limited to low-load resistance exercise where strength and size benefits are limited compared to high-load resistance exercise.
Low-load resistance exercise combined with blood flow restriction (BFR) has been shown in literature to increase both muscle strength and size across different age groups. With BFR, exercise can be performed at substantially lower loads and intensities while generating similar muscular and physiological adaptations seen in high intensity resistance training. For load compromised populations, this reduces the pain during the exercise protocol and leads to overall improvements in physical function.
To provide consistent BFR pressure stimulus to patients, it is recommended to (1) apply a restrictive pressure that is personalized to each individual patient based on the patient’s limb occlusion pressure, and (2) utilize a BFR system that can provide surgical-grade tourniquet autoregulation.
| Technology | Devices | null |
18468102 | https://en.wikipedia.org/wiki/Dead%20leaf%20mantis | Dead leaf mantis | Dead leaf mantis is a common name given to various species of praying mantis that mimic dead leaves. It is most often used in reference to species within genus Deroplatys because of their popularity as exotic pets. Examples include D. desiccata (giant dead leaf mantis), D. lobata (Southeast Asian dead leaf mantis), and D. philippinica (Philippines dead leaf mantis). Other species to which the term may apply include Acanthops falcataria (South American dead leaf mantis), A. falcata (South American dead leaf mantis), and Phyllocrania paradoxa (more common known as the ghost mantis).
| Biology and health sciences | Insects: General | Animals |
18471728 | https://en.wikipedia.org/wiki/Agricultural%20land | Agricultural land | Agricultural land is typically land devoted to agriculture, the systematic and controlled use of other forms of lifeparticularly the rearing of livestock and production of cropsto produce food for humans. It is generally synonymous with both farmland or cropland, as well as pasture or rangeland.
The United Nations Food and Agriculture Organization (FAO) and others following its definitions, however, also use agricultural land or as a term of art, where it means the collection of:
arable land (also known as cropland): here redefined to refer to land producing crops requiring annual replanting or fallowland or pasture used for such crops within any five-year period
permanent cropland: land producing crops which do not require annual replanting
permanent pastures: natural or artificial grasslands and shrublands able to be used for grazing livestock
This sense of "agricultural land" thus includes a great deal of land not devoted to agricultural use. The land actually under annually-replanted crops in any given year is instead said to constitute or "Permanent cropland" includes forested plantations used to harvest coffee, rubber, or fruit but not tree farms or proper forests used for wood or timber. Land able to be used for farming is called . Farmland, meanwhile, is used variously in reference to all agricultural land, to all cultivable land, or just to the newly restricted sense of "arable land". Depending upon its use of artificial irrigation, the FAO's "agricultural land" may be divided into irrigated and non-irrigated land.
In the context of zoning, agricultural land or agriculturally-zoned land refers to plots that are permitted to be used for agricultural activities, without regard to its present use or even suitability. In some areas, agricultural land is protected so that it can be farmed without any threat of development. The Agricultural Land Reserve in British Columbia in Canada, for instance, requires approval from its Agricultural Land Commission before its lands can be removed or subdivided.
Area
Under the FAO's definitions above, agricultural land covers 38.4% of the world's land area as of 2011. Permanent pastures are 68.4% of all agricultural land (26.3% of global land area), arable land (row crops) is 28.4% of all agricultural land (10.9% of global land area), and permanent crops (e.g. vineyards and orchards) are 3.1% (1.2% of global land area).
Total of land used to produce food:
Arable land:
Permanent pastures:
Permanent crops:
In 2022, the global agricultural land area was 4.78 billion hectares (ha), down from 4.79 billion hectares in 2021. One-third of the total agricultural land was cropland (1.58 billion ha in 2021), which increased by 6 percent (0.09 billion ha).
Asia had the largest share of the global cropland area in 2021 (37 percent), followedby the Americas (24 percent), Africa (19 percent), Europe (18 percent) and Oceania (2 percent). There were differences in cropland expansion in the different regions during this period – Oceania and Africa both had rapid growth in cropland area (33 percent and 27 percent), while Asia and the Americas had more moderate growth (4 percent and 2 percent). The cropland area of Europe declined between 2000 and 2021 by 5 percent. As aresult, the cropland area of Africa overtook that of Europe in 2018.
Approximately 30 percent of global cropland and permanent meadows and pastures can be found in three countries. In 2021, 12 percent of global permanent meadows and pastures belonged to China, 10 percent to Australia, and 8 percent to the United States of America. For the same year, the largest share of global cropland was in India (11 percent), followed by the United States of America (10 percent) and China (8 percent).
Cropland area per capita decreased in all regions between 2000 and 2021 as population increased faster than the cropland area. The world average declined by 18 percent to 0.20 ha per capitain 2021; the decrease was the largest in Africa (−25 percent, to0.21 ha per capita), followed by the Americas and Asia (−17 percent each,to 0.37 ha per capita and 0.13 ha per capita, respectively), Europe and Oceania (−7 percent each, to 0.39 haper capita and 0.77 ha per capita, respectively). The countries with the highest croplandarea per capita are Kazakhstan, Australia and Canada, due to vast areas of land available.
Globally, the total amount of permanent pasture according to the FAO has been in decline since 1998, in part due to a decrease of wool production in favor of synthetic fibers (such as polyester) and cotton.
The decrease of permanent pasture, however, does not account for gross conversion (e.g. land extensively cleared for agriculture in some areas, while converted from agriculture to other uses elsewhere) and more detailed analyses have demonstrated this. For example, Lark et al. 2015 found that in the United States cropland increased by 2.98 million acres from 2008 to 2012 (comprising converted to agriculture, and converted from agriculture).
Source: Helgi Library, World Bank, FAOSTAT
Agricultural land market
Prices and rents for agricultural land depend on supply and demand.
Prices/rents rise when the supply of farmland on the market reduces. Landholders then put more land on the market – causing prices to fall. Conversely, land prices/rents fall when the demand for agricultural land declines because of falls in the returns from holding and using it. The immediate triggers for falls in land demand might be reductions in the demand for farm produce or in relevant government subsidies and tax reliefs.
Russia
The cost of Russian farmland is as little as €1,500–2,000 (£1,260–1,680) per hectare (ha) (£1,260–1,680). This is comparatively inexpensive. Poor-quality farmland in France and Spain is sold at no lower than €10,000/ha.
The average Russian farm measures 150 hectares (370 acres). The most prevalent crops in Russia are wheat, barley, corn, rice, sugar beet, soy beans, sunflower, potatoes and vegetables. Russian farmers harvested roughly 85–90 million tonnes of wheat annually in the years around 2010. Russia exported most to Egypt, Turkey and Iran in 2012; China was a significant export market as well. The average yield from the Krasnodar region was between 4 and 5 tonnes per ha, while the Russian average was only 2t/ha. The Basic Element Group, a conglomerate owned by Oleg Deripaska, is one of Russia's leading agricultural producers, and owns or manages 109,000ha of Russian farmland, out of 90m actual and 115m total (0.12% actual).
Ukraine
In 2013, Ukraine was ranked third in corn production and sixth in wheat production. It was the main supplier of corn, wheat, and rape to Europe, although it is unclear whether the internal supply from countries like France were accounted in this calculation. Ukrainian farmers achieve 60% of the output per unit area of their North American competitors. UkrLandFarming PLC produces, from 650,000 hectares (1.6m acres), corn, wheat, barley, sugar beet, and sunflowers. Until 2014, the chief Ukrainian export terminal was the Crimean port of Sevastopol.
United States
Prime farmland in Illinois is valued, as of August 2018, at $26,000 a hectare. Average cropland value in the Midwest according to 2020 data from the US Department of Agriculture is $4,607 per acre (about $11,000 per hectare).
| Technology | Basics_2 | null |
3476955 | https://en.wikipedia.org/wiki/Shear%20%28geology%29 | Shear (geology) | In geology, shear is the response of a rock to deformation usually by compressive stress and forms particular textures. Shear can be homogeneous or non-homogeneous, and may be pure shear or simple shear. Study of geological shear is related to the study of structural geology, rock microstructure or rock texture and fault mechanics.
The process of shearing occurs within brittle, brittle-ductile, and ductile rocks. Within purely brittle rocks, compressive stress results in fracturing and simple faulting.
Rocks
Rocks typical of shear zones include mylonite, cataclasite, S-tectonite and L-tectonite, pseudotachylite, certain breccias and highly foliated versions of the wall rocks.
Shear zone
A shear zone is a tabular to sheetlike, planar or curviplanar zone composed of rocks that are more highly strained than rocks adjacent to the zone. Typically this is a type of fault, but it may be difficult to place a distinct fault plane into the shear zone. Shear zones may form zones of much more intense foliation, deformation, and folding. En echelon veins or fractures may be observed within shear zones.
Many shear zones host ore deposits as they are a focus for hydrothermal flow through orogenic belts. They may often show some form of retrograde metamorphism from a peak metamorphic assemblage and are commonly metasomatised.
Shear zones can be only inches wide, or up to several kilometres wide. Often, due to their structural control and presence at the edges of tectonic blocks, shear zones are mappable units and form important discontinuities to separate terranes. As such, many large and long shear zones are named, identical to fault systems.
When the horizontal displacement of this faulting can be measured in the tens or hundreds of kilometers of length, the fault is referred to as a megashear. Megashears often indicate the edges of ancient tectonic plates.
Mechanisms of shearing
The mechanisms of shearing depend on the pressure and temperature of the rock and on the rate of shear which the rock is subjected to. The response of the rock to these conditions determines how it accommodates the deformation.
Shear zones which occur in more brittle rheological conditions (cooler, less confining pressure) or at high rates of strain, tend to fail by brittle failure; breaking of minerals, which are ground up into a breccia with a milled texture.
Shear zones which occur under brittle-ductile conditions can accommodate much deformation by enacting a series of mechanisms which rely less on fracture of the rock and occur within the minerals and the mineral lattices themselves. Shear zones accommodate compressive stress by movement on foliation planes.
Shearing at ductile conditions may occur by fracturing of minerals and growth of sub-grain boundaries, as well as by lattice glide. This occurs particularly on platy minerals, especially micas.
Mylonites are essentially ductile shear zones.
Microstructures of shear zones
During the initiation of shearing, a penetrative planar foliation is first formed within the rock mass. This manifests as realignment of textural features, growth and realignment of micas and growth of new minerals.
The incipient shear foliation typically forms normal to the direction of principal shortening, and is diagnostic of the direction of shortening. In symmetric shortening, objects flatten on this shear foliation much the same way that a round ball of treacle flattens with gravity.
Within asymmetric shear zones, the behavior of an object undergoing shortening is analogous to the ball of treacle being smeared as it flattens, generally into an ellipse. Within shear zones with pronounced displacements a shear foliation may form at a shallow angle to the gross plane of the shear zone. This foliation ideally manifests as a sinusoidal set of foliations formed at a shallow angle to the main shear foliation, and which curve into the main shear foliation. Such rocks are known as L-S tectonites.
If the rock mass begins to undergo large degrees of lateral movement, the strain ellipse lengthens into a cigar shaped volume. At this point shear foliations begin to break down into a rodding lineation or a stretch lineation. Such rocks are known as L-tectonites.
Ductile shear microstructures
Very distinctive textures form as a consequence of ductile shear. An important group of microstructures observed in ductile shear zones are S-planes, C-planes and C' planes.
S-planes or schistosité planes are generally defined by a planar fabric caused by the alignment of micas or platy minerals. Define the flattened long-axis of the strain ellipse.
C-planes or cisaillement planes form parallel to the shear zone boundary. The angle between the C and S planes is always acute, and defines the shear sense. Generally, the lower the C-S angle the greater the strain.
The C' planes, also known as shear bands and secondary shear fabrics, are commonly observed in strongly foliated mylonites especially phyllonites, and form at an angle of about 20 degrees to the S-plane.
The sense of shear shown by both S-C and S-C' structures matches that of the shear zone in which they are found.
Other microstructures which can give sense of shear include:
sigmoidal veins
mica fish
rotated porphyroclasts
asymmetric boudins (Figure 1)
asymmetric folds
Transpression
Transpression regimes are formed during oblique collision of tectonic plates and during non-orthogonal subduction. Typically a mixture of oblique-slip thrust faults and strike-slip or transform faults are formed. Microstructural evidence of transpressional regimes can be rodding lineations, mylonites, augen-structured gneisses, mica fish and so on.
A typical example of a transpression regime is the Alpine Fault zone of New Zealand, where the oblique subduction of the Pacific Plate under the Indo-Australian Plate is converted to oblique strike-slip movement. Here, the orogenic belt attains a trapezoidal shape dominated by oblique splay faults, steeply-dipping recumbent nappes and fault-bend folds.
The Alpine Schist of New Zealand is characterised by heavily crenulated and sheared phyllite. It is being pushed up at the rate of 8 to 10 mm per year, and the area is prone to large earthquakes with a south block up and west oblique sense of movement.
Transtension
Transtension regimes are oblique tensional environments. Oblique, normal geologic fault and detachment faults in rift zones are the typical structural manifestations of transtension conditions. Microstructural evidence of transtension includes rodding or stretching lineations, stretched porphyroblasts, mylonites, etc.
| Physical sciences | Structural geology | Earth science |
3479667 | https://en.wikipedia.org/wiki/Parton%20%28particle%20physics%29 | Parton (particle physics) | In particle physics, the parton model is a model of hadrons, such as protons and neutrons, proposed by Richard Feynman. It is useful for interpreting the cascades of radiation (a parton shower) produced from quantum chromodynamics (QCD) processes and interactions in high-energy particle collisions.
History
The parton model was proposed by Richard Feynman in 1969, used originally for analysis of high-energy hadron collisions. It was applied to electron-proton deep inelastic scattering by James Bjorken and Emmanuel Anthony Paschos. Later, with the experimental observation of Bjorken scaling, the validation of the quark model, and the confirmation of asymptotic freedom in quantum chromodynamics, partons were matched to quarks and gluons. The parton model remains a justifiable approximation at high energies, and others have extended the theory over the years.
Murray Gell-Mann preferred to use the term "put-ons" to refer to partons.
In 1994, partons were used by Leonard Susskind to model holography.
Model
Any hadron (for example, a proton) can be considered as a composition of a number of point-like constituents, termed "partons".
Component particles
Just as accelerated electric charges emit QED radiation (photons), the accelerated coloured partons will emit QCD radiation in the form of gluons. Unlike the uncharged photons, the gluons themselves carry colour charges and can therefore emit further radiation, leading to parton showers.
Reference frame
The hadron is defined in a reference frame where it has infinite momentum – a valid approximation at high energies. Thus, parton motion is slowed by time dilation, and the hadron charge distribution is Lorentz-contracted, so incoming particles will be scattered "instantaneously and incoherently".
Partons are defined with respect to a physical scale (as probed by the inverse of the momentum transfer). For instance, a quark parton at one length scale can turn out to be a superposition of a quark parton state with a quark parton and a gluon parton state together with other states with more partons at a smaller length scale. Similarly, a gluon parton at one scale can resolve into a superposition of a gluon parton state, a gluon parton and quark-antiquark partons state and other multiparton states. Because of this, the number of partons in a hadron actually goes up with momentum transfer. At low energies (i.e. large length scales), a baryon contains three valence partons (quarks) and a meson contains two valence partons (a quark and an antiquark parton). At higher energies, however, observations show sea partons (nonvalence partons) in addition to valence partons.
Parton distribution functions
A parton distribution function (PDF) within so called collinear factorization is defined as the probability density for finding a particle with a certain longitudinal momentum fraction x at resolution scale Q2. Because of the inherent non-perturbative nature of partons which cannot be observed as free particles, parton densities cannot be calculated using perturbative QCD. Within QCD one can, however, study variation of parton density with resolution scale provided by external probe. Such a scale is for instance provided by a virtual photon with virtuality Q2 or by a jet. The scale can be calculated from the energy and the momentum of the virtual photon or jet; the larger the momentum and energy, the smaller the resolution scale—this is a consequence of Heisenberg's uncertainty principle. The variation of parton density with resolution scale has been found to agree well with experiment; this is an important test of QCD.
Parton distribution functions are obtained by fitting observables to experimental data; they cannot be calculated using perturbative QCD. Recently, it has been found that they can be calculated directly in lattice QCD using large-momentum effective field theory.
Experimentally determined parton distribution functions are available from various groups worldwide. The major unpolarized data sets are:
ABM by S. Alekhin, J. Bluemlein, S. Moch
CTEQ, from the CTEQ Collaboration
GRV/GJR, from M. Glück, P. Jimenez-Delgado, E. Reya, and A. Vogt
HERA PDFs, by H1 and ZEUS collaborations from the Deutsches Elektronen-Synchrotron center (DESY) in Germany
MSHT/MRST/MSTW/MMHT, from A. D. Martin, R. G. Roberts, W. J. Stirling, R. S. Thorne, and collaborators
NNPDF, from the NNPDF Collaboration
The LHAPDF library provides a unified and easy-to-use Fortran/C++ interface to all major PDF sets.
Generalized parton distributions (GPDs) are a more recent approach to better understand hadron structure by representing the parton distributions as functions of more variables, such as the transverse momentum and spin of the parton. They can be used to study the spin structure of the proton, in particular, the Ji sum rule relates the integral of GPDs to angular momentum carried by quarks and gluons. Early names included "non-forward", "non-diagonal" or "skewed" parton distributions. They are accessed through a new class of exclusive processes for which all particles are detected in the final state, such as the deeply virtual Compton scattering. Ordinary parton distribution functions are recovered by setting to zero (forward limit) the extra variables in the generalized parton distributions. Other rules show that the electric form factor, the magnetic form factor, or even the form factors associated to the energy-momentum tensor are also included in the GPDs. A full 3-dimensional image of partons inside hadrons can also be obtained from GPDs.
Simulation
Parton showers simulations are of use in computational particle physics either in automatic calculation of particle interaction or decay or event generators, in order to calibrate and interpret (and thus understand) processes in collider experiments. They are particularly important in large hadron collider (LHC) phenomenology, where they are usually explored using Monte Carlo simulation.
The scale at which partons are given to hadronization is fixed by the Shower Monte Carlo program. Common choices of Shower Monte Carlo are PYTHIA and HERWIG.
| Physical sciences | Subatomic particles: General | Physics |
3482079 | https://en.wikipedia.org/wiki/Sodium%20orthovanadate | Sodium orthovanadate | Sodium orthovanadate is the inorganic compound with the chemical formula . It forms a dihydrate . Sodium orthovanadate is a salt of the oxyanion. It is a colorless, water-soluble solid.
Synthesis and structure
Sodium orthovanadate is produced by dissolving vanadium(V) oxide in a solution of sodium hydroxide:
The salt features tetrahedral anion centers linked to octahedral cation sites.
Condensation equilibria
Like many oxometalates, orthovanadate is subject to a number of reactions, which have been analyzed by 51V NMR studies. At high pH, ions exist in equilibrium with . At lower pH's, condensation ensues to give various polyoxovanadates. Ultimately, decavanadate is formed.
Biochemistry
Vanadates exhibit a variety of biological activities, in part because they serve as structural mimics of phosphates. It acts as a competitive inhibitor of ATPases, alkaline and acid phosphatases, and protein-phosphotyrosine phosphatases, and its inhibitory effects can be reversed by dilution or the addition of ethylenediaminetetraacetic acid (EDTA).
Orthovanadate is activated by boiling and adjusting pH to ~10; this depolymerizes decavanadate into the active inhibitor, monovanadate.
| Physical sciences | Metallic oxyanions | Chemistry |
3482374 | https://en.wikipedia.org/wiki/Military%20glider | Military glider | Military gliders (an offshoot of common gliders) have been used by the militaries of various countries for carrying troops (glider infantry) and heavy equipment to a combat zone, mainly during the Second World War. These engineless aircraft were towed into the air and most of the way to their target by military transport planes, e.g., C-47 Skytrain or Dakota, or bombers relegated to secondary activities, e.g., Short Stirling. Most military gliders do not soar, although there were attempts to build military sailplanes as well, such as the DFS 228.
Once released from the tow craft near the front, they were to land on any convenient open terrain close to the target, hopefully with as little damage to the cargo and crew as possible, as most landing zones (LZ) were far from ideal. The one-way nature of the missions meant that they were treated as semi-expendable leading to construction from common and inexpensive materials such as wood. Most nations seriously attempted to recover as many as possible, to re-use them, so they were not originally intended to be disposable, although resource-rich nations like the US sometimes used them as if they were, since it was easier than recovering them.
Troops landing by glider were referred to as air-landing as opposed to paratroops. Landing by parachute caused the troops to be spread over a large drop-zone and separated from other airdropped equipment, such as vehicles and anti-tank guns. Gliders, on the other hand, could land troops and ancillaries in greater concentrations precisely at the target landing area. Furthermore, the glider, once released at some distance from the actual target, was effectively silent and difficult for the enemy to identify. Larger gliders were developed to land heavy equipment like anti-tank guns, anti-aircraft guns, small vehicles, such as jeeps, and also light tanks (e.g., the Tetrarch tank). This heavier equipment made otherwise lightly armed paratroop forces a much more capable force. The Soviets also experimented with ways to deliver light tanks by air, including the Antonov A-40, a gliding tank with detachable wings.
By the time of the Korean War, helicopters had largely replaced gliders. Helicopters have the advantage of being able to extract soldiers, in addition to delivering them to the battlefield with more precision. Also, advances in powered transport aircraft had been made, to the extent that even light tanks could be dropped by parachute. And after the widespread use of radar in the military, silence in the air is no longer sufficient for concealment.
Development
The development of modern gliders was spurred by the Versailles Treaty following World War I, under the terms of which Germany was prohibited from constructing certain high powered airplanes. As a result, German aircraft designers turned their attention toward the practical development of unpowered aircraft, with a pilot remaining in the air in a glider for more than 20 minutes and a national glider competition emerging by 1922.
The early sporting objectives of gliders were quickly overtaken in the Soviet Union and in Germany by military applications, mainly the training of pilots. By 1934, the Soviet Union had ten gliding schools and 57,000 glider pilots had gained licences.
In 1932, the Soviet Union demonstrated the TsK Komsula, a four-place glider, designed by GF Groschev that could also be used for cargo. Larger gliders were then developed culminating in an 18-seater at the military institute in Leningrad in 1935. Luftwaffe Colonel Kurt Student visited Moscow as part of the military collaboration programme with the Soviet Union. He reported back to his superiors in Berlin details of a 1,500 man parachute drop and the large transport gliders that he had seen. The Luftwaffe opened a parachute school as a result in 1937. Further field testing convinced Student that a vehicle was needed to deliver the heavy weapons for the lightly armed parachute troops. This idea was dismissed until October 1938 by which time Student had risen to major-general and was appointed Inspector of Airborne Forces. Development of a troop-carrying glider was assigned to Hans Jacobs of the Deutsche Forschungsanstalt für Segelflug to develop the DFS 230 which could carry 9–10 fully equipped troops or 1,200 kg (2,800 pounds).
German military glider
The Germans were the first to use gliders in warfare, most famously during the assault of the Eben Emael fortress and the capture of the bridges over the Albert Canal at Veldwezelt, Vroenhoven and Kanne on May 10, 1940, in which 41 DFS 230 gliders carrying 10 soldiers each were launched behind Junkers Ju 52s. Ten gliders landed on the grassed roof of the fortress. Only twenty minutes after landing the force had neutralized the fortress at a cost of six dead and twenty wounded. Hitler was anxious to gain maximum publicity and so several foreign attachés were given guided tours of the fortress. Consequently, the British, American and Japanese became quickly aware of the methods that had been used. By mid-1940, both Japan and Britain had active glider programs.
Development then began of even larger gliders such as the Gotha Go 242 (23 trooper) and Messerschmitt Me 321 (130 trooper) to transport heavy armaments in anticipation of Operation Sea Lion and Operation Barbarossa.
Gliders were also used by Germany in Greece in 1941. On April 26, 1941, the troops from six DFS 230 gliders captured the bridge over the Corinth Canal accompanied by 40 plane-loads of German paratroopers. (Fortuitously, the British were able to demolish the bridge a few hours later.) Next, General Student then convinced Hitler that Crete could be captured using only airborne troops. Consequently, on May 20, 1941, 500 German transport aircraft carrying paratroopers and 74 DFS 230 gliders took off from the Greek mainland. During the capture of the island, 5,140 German airborne troops were either killed or wounded out of the 13,000 sent. Among the 350 German planes destroyed in the operation, half had been Ju 52s, which seriously depleted the force needed for the invasion of the Soviet Union shortly after. As a result, Hitler vowed never to use his airborne force in such large numbers again.
Some German glider operations continued later in the war, some examples being the rescue operation of Benito Mussolini at Gran Sasso and emergency re-supply operations in Russia, North Africa and Eastern Europe towards the end of the war. The Junkers Ju 322 Mammut ("Mammoth") was the largest such glider ever built, but it was never used operationally. Not all military gliders were planned for transport. The Blohm & Voss BV 40 was a German glider fighter designed to attack Allied bomber formations but was not used.
British military gliders
The British glider development started in mid-1940, prompted by the assault on Eben Emael. Among the types developed were the 28 trooper Airspeed Horsa and the 7-ton capacity General Aircraft Hamilcar cargo glider. The Hamilcar could carry vehicles, anti-tank guns and light tanks into action. The General Aircraft Hotspur – originally planned as a compact assault glider carrying a small number of troops – was used for training the British Army pilots who formed the Glider Pilot Regiment. The Slingsby Hengist was a backup design which was not required when the similar capacity American-built Waco CG-4 (given the British service name "Hadrian") became available in large numbers through lend-lease. Four hundred of the 3,600 Horsas built were supplied to the USAAF.
The most famous British actions using gliders included the unsuccessful Operation Freshman, against a German heavy water plant in Norway in 1942; and the capture of the Caen canal and Orne river bridges in a coup-de-main operation at the very start of the invasion of Normandy. Other glider actions included Operation Dragoon (the invasion of southern France), Operation Market Garden (the landing at Arnhem Bridge to try and seize a bridgehead over the lower Rhine) and Operation Varsity (crossing of the Rhine). Out of the 2,596 gliders dispatched for Operation Market Garden, 2,239 were effective in delivering men and equipment to their designated landing zones.
Although gliders are still used in the Royal Air Force in the Royal Air Force Gliding & Soaring Association and for cadet training by the Air Training Corps, they are not used in combat operations. No troop-carrying gliders have been in British service since 1957.
American military gliders
United States Army, Army Air Forces, and Air Force
Major General Henry "Hap" Arnold, Acting Deputy Chief of Staff for Air (becoming Commanding General of the United States Army Air Forces on March 9, 1942), initiated a study with view to develop a glider capable of being towed by aircraft. This directive was set into motion through Classified Technical Instructions (CTI-198 on 24 February 1941, and CTI-203 on 4 March 1941), which authorized the procurement of 2-, 8-, and 15-place gliders and equipment. Eleven companies were invited to participate in the experimental glider program, but only four responded with any interest, Frankfort Sailplane Company (XCG-1, XCG-2), Waco Aircraft Company (XCG-3, XCG-4), St. Louis Aircraft Corp. (XCG-5, XCG-6), and Bowlus Sailplanes (XCG-7, XCG-8). Only Waco Aircraft Company was able to deliver the experimental glider prototypes that satisfied the requirements of Materiel Command, the eight-seat Waco CG-3 (modified to become a production nine-seat glider) and the fifteen-seat Waco CG-4. In October 1941, Lewin B. Barringer was made Glider Specialist, Air Staff, HQ of the Army Air Forces, answering to General Arnold, and placed in charge of the glider program. The shock of the Japanese attack on Pearl Harbor on 7 December 1941 prompted the United States to set the number of glider pilots needed at 1,000 to fly 500 eight-seat gliders and 500 fifteen-seat gliders. The number of pilots required was increased to 6,000 by June 1942. After Barringer was lost at sea on a flight to Africa in January 1943, the program came under direction of Richard C. du Pont. Bigger gliders, such as the 30-troop Waco CG-13A and the 42-troop Laister-Kauffman CG-10A were designed later.
The most widely used type was the Waco CG-4A, which was first used in the invasion of Sicily in July 1943 and participated in the D-Day assault on France on 6 June 1944, and in other important airborne operations in Europe, including Operation Market Garden in September 1944 and the crossing the Rhine in March 1945, and in the China-Burma-India Theater. The CG-4A was constructed of a metal and wood frame covered with fabric, manned by a crew of two and with an allowable normal cargo load of 3,710 lb, allowing it to carry 13 combat-equipped troops or a jeep or small artillery piece. The CG-10 could hold 10,850 lb of cargo, such as two howitzers, at a time. The final glider mission of the war was at Luzon on 23 June 1945. By the end of the war, the United States had built 14,612 gliders of all types and had trained over 6,000 glider pilots. The designs of the Waco Aircraft Company were also produced by a wide variety of manufacturers including Ford Motor Company and Cessna Aircraft Company as well as furniture, piano and coffin manufacturers.
Following World War II, the United States maintained only one regiment of gliders. Gliders were used in military exercises in 1949, but glider operations were deleted from the United States Army's capabilities on 1 January 1953. However, the United States Air Force continues to use sailplanes at the United States Air Force Academy to train cadets in the fundamentals of flight.
United States Navy and Marine Corps
In April 1941, United States Navy officer Marc Mitscher proposed that the Navy develop amphibious gliders with flying-boat hulls with a goal of deploying an amphibious glider force capable of delivering an entire United States Marine Corps brigade of 715 men to a hostile beachhead, the gliders to be towed by Consolidated PBY-5A Catalina amphibian aircraft. The Navy's Bureau of Aeronautics developed specifications for two types of amphibious glider, a single-hulled type which could carry 12 passengers and a twin-hulled type that could carry 24 passengers. Two companies, the Allied Aviation Corporation and the Bristol Aeronautical Corporation, received contracts to produce 100 gliders, and plans called for the procurement of 12,000 more amphibious gliders if the concept proved successful.
No twin-hulled glider was built, but each company constructed the prototype of a single-hulled amphibious glider, the XLRA-1 by Allied Aviation and the XLRQ-1 by Bristol Aeronautical. The two prototypes made their first flights in early 1943, but by the time they did the Navy and Marine Corps already had concluded that the use of gliders to deliver Marines to beachheads was impractical. No further examples of the two glider types were built, and the Navy officially terminated the amphibious glider program on 27 September 1943. Testing of the two prototypes continued until early December 1943, apparently in connection with the development of a glider bomb.
The Marine Corps established a glider training unit in early 1942 at Marine Corps Recruit Depot Parris Island, South Carolina, using non-amphibious Pratt-Read LNE-1 and Schweizer LNS-1 gliders. In addition, the Navy took delivery during World War II of 15 U.S. Army Air Forces Waco CG-4A non-amphibious gliders for evaluation under the Navy designation LRW-1. Neither of these initiatives resulted in operational use of gliders by the U.S. Navy or Marine Corps.
Soviet military gliders
The Soviet Union built the world's first military gliders starting in 1932, including the 16-seat Grokhovski G63, though no glider was built in quantity until World War II. During the war, there were only two light gliders built in series: Antonov A-7 and Gribovski G-11 – about 1,000 altogether. A medium glider, the KC-20, was built in a small series. They were used mostly for providing partisans in Belarus with supplies and armament in 1942–1943. On 21 September 1943, 35 gliders were used in the Dnepr crossing. Later, other types of gliders were built: the Cybin C-25 (25 trooper) in 1944, the Yakovlev Yak-14 (35 trooper) in 1948, and the Ilyushin Il-32 (60 trooper) also in 1948. In 1950, a Yak-14 became the first glider to fly over the North Pole.
The Soviet Union maintained three glider infantry regiments until 1965. However, Soviet Air Force transport gliders were gradually withdrawn from service with the arrival of turboprop transports like the Antonov An-12 and Antonov An-24, which entered service in the late 1950s.
| Technology | Military aviation | null |
757240 | https://en.wikipedia.org/wiki/Plastic%20mulch | Plastic mulch | Plastic mulch is a product used in plasticulture in a similar fashion to mulch, to suppress weeds and conserve water in crop production and landscaping. Certain plastic mulches also act as a barrier to keep methyl bromide, both a powerful fumigant and ozone depleter, in the soil. Crops grow through slits or holes in thin plastic sheeting. Plastic mulch is often used in conjunction with drip irrigation. Some research has been done using different colors of mulch to affect crop growth. Use of plastic mulch is predominant in large-scale vegetable growing, with millions of acres cultivated under plastic mulch worldwide each year.
Disposal of plastic mulch is an environmental problem. Technologies exist to provide for the recycling of used/disposed plastic mulch into viable plastic resins for re-use in the plastics manufacturing industry. However these methods are not very effective due to contamination by agrochemicals of the plastic. Other concerns include residual microplastics in the soil which can have negative effects on soil ecologies, including microbes and earthworms.
History
The idea of using polyethylene film as mulch in plant production saw its beginnings in the mid-1950s. Dr. Emery M. Emmert of the University of Kentucky was one of the first to recognize the benefits of using LDPE (low-density polyethylene) and HDPE (high-density polyethylene) film as mulch in vegetable production. Emmert also wrote on other topics such as the use of plastic for greenhouses instead of glass and plastic in field high tunnels. Approximately of agricultural land utilize polyethylene mulch and similar row covers for crop production in the world.
Laying plastic polythene (mulch) down over mounds formed in the soil was also pioneered in New Zealand in the mid fifties by strawberry growers in the Auckland area.
By 1960-61 all strawberries grown commercially in New Zealand were grown through black polythene usually laid by hand. The plastic promoted growth, conserved moisture brought on early fruiting and restricted weed infestation.
The earliest polythene laying machines were in use in New Zealand by the mid 1960s and were very similar to the machines sold today. The very first machines were designed by growers and built by small engineering/fabrication workshops, usually under the careful guidance and supervision of the farmer. Each machine for many years was generally similar to the last, with the occasional modification to improve performance.
Benefits
The use of plastic mulches along with the use of drip irrigation has many benefits such as:
Soil temperature
The use of plastic mulch alters soil temperature. Dark mulches and clear mulches applied to the soil intercept sunlight and warm the soil, allowing earlier planting as well as encouraging faster growth early in the growing season. White mulch reflects heat from the sun, effectively reducing soil temperature. This reduction in temperature may help establish plants in mid-summer when cooler soil might be required.
Soil moisture retention
Plastic mulches reduce the amount of water lost from the soil due to evaporation. This means less water will be needed for irrigation. Plastic mulches also aid in evenly distributing moisture to the soil, which reduces plant stress.
Weed management
Plastic mulches prevent sunlight from reaching the soil which can inhibit most annual and perennial weeds. Clear plastics prevent weed growth. Holes in the mulch for plants tend to be the only pathway for weeds to grow.
Reduction in the leaching of fertilizer
The use of drip irrigation in conjunction with plastic mulch allows one to reduce leaching of fertilizers. Using drip irrigation eliminates the use of flood and furrow irrigation that applies large quantities of water to the soil, which in turn tends to leach nitrogen and other nutrients to depths below the root zone. Drip irrigation applies lower amounts of water with fertilizers injected and thus these fertilizers are applied to the root zone as needed. This also reduces the amount of fertilizer needed for adequate plant growth when compared to broadcast fertilization.
Improved crop quality
Plastic mulches keep ripening fruits off of the soil. This reduced contact with the soil decreases fruit rot as well as keeps the fruit and vegetables clean. This is beneficial for the production of strawberries, for example.
Reduction in soil compaction
The plastic mulch covering the soil decreases the crusting effect of rain and sunlight. The reduction in weed quantity means a decreased need for mechanical cultivation. Weed control between beds of plastic can be done using directly applied herbicides and through mechanical means. The soil underneath the plastic mulch stays loose and well aerated, with the mulch protecting the soil it covers from erosion.
Reduction in root damage
The use of plastic mulch creates a practically weed-free area around the plant, removing the need for cultivation except between the rows of plastic. Root damage associated with cultivation is therefore eliminated. Due to these factors, the use of plastic mulch can lead to an improvement in the overall growth of the plant.
Disadvantages
There are many disadvantages to using plastic mulches in crop production as well.
Cost
The benefits from using plastic mulch come at a higher cost than planting in bare soil. These costs include equipment, the plastic film used as the mulch, transplanters designed for plastic beds, and additional labor during installation and removal of mulch films. Specialized mulch application equipment must be used to install plastic mulch beds into a field. These machines shape the soil and apply the plastic to the prepared soil. Transplanters designed for plastic mulch can be used to plant the desired crop. Hand transplanting is an option, but this is rather inefficient. The removal of plastic mulch also contributes to a higher cost through additional labor and equipment needed. Specialized designed undercutting equipment can be used to remove the plastic from the field after harvest.
Environmental concerns
If conventional plastics (e.g. PE) are used as mulch films, they are likely to accumulate in soil, since the removal and the correct disposal of these plastics are technically and economically burdensome. This accumulation could cause both crop yield reduction and environmental problems.
Biodegradable polymers are polymers that can be degraded by the naturally occurring microbial community in an environmental system. They provide a more sustainable alternative to conventionally used plastics for mulch films. Providing the same benefits as detailed above, the problem of plastic accumulation in soils could be solved. Aliphatic polyesters and aliphatic-aromatic co-polyesters have shown to be promising groups of biodegradable polymers.
Application
The use of plastic mulch requires a unique application process to ensure proper placement of the plastic film. This application process begins with preparing the field the same way one would for a flat seed bed. The bed must be free of large soil clods and organic residue. A machine called a plastic layer or a bed shaper is pulled over the field creating a row of plastic mulch covering a planting bed. These beds can be a flat bed which simply means the surface of the plastic mulch is level with the inter-row soil surface. Machines that form raised beds create a plastic surface higher than the inter-row soil surface. The basic concept of the plastic bed shaper is a shaping box which creates the bed, that is then covered by plastic via a roller and two coulters that cover the edges of the plastic film to hold the plastic the soil's surface. These plastic layers also place the drip irrigation line under the plastic while the machine lays the plastic. It is somewhat important that the plastic is rather tight. This becomes important in the planting process.
Planting
Planting also requires specialized planting equipment. The most common planting equipment is a waterwheel type transplanter. The waterwheel transplanter utilizes a rotating drum or drums with spikes at set intervals. The drum or drums have a water supply that continuously fills the drum with water. The transplanter rolls the spiked drum over the bed of plastic. As the drum presses a spike into the plastic, a hole is punched and water flows into the punched hole. A rider on the transplanter can then place a plant in the hole. These drums can have multiple rows and varied intervals to create the desired spacing for that particular crop.
| Technology | Soil and soil management | null |
757361 | https://en.wikipedia.org/wiki/EuroVelo | EuroVelo | EuroVelo is a network of 17 long-distance cycling routes criss-crossing Europe, with 2 more in early construction across various stages of completion. When completed, the EuroVelo network's total length will be almost . more than were in place. EuroVelo is a project of the European Cyclists' Federation (ECF). The multinational project aims to connect 40 countries via the 19 unique routes across the European continent.
EuroVelo routes can be used for bicycle touring across the continent, as well as by local people making short journeys. The routes are made of both existing national bike routes — such as the Dutch LF-Routes, the German D-Routes, the French véloroute "SN3V" and the British National Cycle Network — and existing general purpose roads, together with new stretches of cycle routes to connect them.
History
The idea of creating a network of international cycle routes spanning Europe started in 1995. It was initially coordinated by the ECF, De Frie Fugle (Denmark) and Sustrans (UK) and the original plan was to create 12 long-distance cycling routes.
Since August 2007, the ECF has assumed full responsibility for the project. Despite sometimes tight financial constraints, the EuroVelo project has already begun to fulfil the vision of its founders with sections of the network being implemented in countries as far apart as Finland, Cyprus, Spain and the UK. In addition, the EuroVelo brand has become widely known.
There have been various changes to the network over the years, most notably the addition of two new routes — EuroVelo 13 (the Iron Curtain Trail) and EuroVelo 15 (the Rhine Cycle Route) — in September 2011, which are the longest and shortest of the EuroVelo routes.
Future expansion
In September 2023, the ECF announced that the Iberian Cycle route connecting Lisbon with Pamplona via Madrid is set to become the future EuroVelo 16 route by 2028 with a length of 1,896 km.
Main points on the EuroVelo routes
Routes EV10 and EV12 are a circular tour
Connections to other EV routes are in parentheses
Odd routes are heading north–south, even routes are heading west–east
Route information
EuroVelo 1 – Atlantic Coast Route
Stretching the length of the continent, from North Cape, Norway to Valença, Portugal, the EV1 connects Norway, Scotland, Northern Ireland, the Republic of Ireland, Wales, the West Country of England, France, Spain and Portugal.
EuroVelo 2 – Capitals Route
EV2 runs between Galway, Ireland to Moscow, Russia visiting some capital cities along the way, from Eyre Square to Red Square.
Between The Hague in the Netherlands and the German-Polish border, the EV2 follows the bicycle route called European Bicycle Route R1 or Euro-Route R1, an international long-distance cycling route connecting Boulogne-sur-Mer in France with St Petersburg in Russia.
EuroVelo 3 – Pilgrims Route
EV3 goes from Trondheim in Norway to Santiago de Compostela in Spain. The route follows traces of old roads used for pilgrimages in the Middle Ages. The route passes through Norway, Sweden, Denmark, Germany, Belgium, France and Spain. Most of these countries have a developed network of bicycle routes used as part of the EV3.
EuroVelo 4 – Central Europe Route
The EV4 goes from Roscoff, France to Kyiv, Ukraine, going through France, Belgium, The Netherlands, Germany, Czechia, Poland, and Ukraine.
EuroVelo 5 – Via Romea Francigena
The EV5 route is inspired by the Via Francigena, a pilgrimage route from London to Rome first recorded by Archbishop of Canterbury Sigeric in the 10th century AD. However, the route of the true Via Francigena is an almost straight line path from London to Rome, while the EuroVelo 5 route takes a more easterly route that passes through Brussels, Luxembourg and Strasbourg in the Alsace. It then follows the Franco-German border, passes through Switzerland following Swiss National Bike Route no. 3, before crossing the Alps at the Gotthard Pass. It then passes through Italy (more closely following Sigeric's route) to Rome before continuing on to the Adriatic port city of Brindisi.
EuroVelo 6 – River Route
Running from Saint-Nazaire on the mouth of the river Loire along that river eastward through France, EV6 passes over the border to Switzerland to Lake Constance and then on to Tuttlingen in Germany, where it begins its way down the Danube following the Donauradweg (Danube Cycle Route). It follows that river, Europe's second longest, through Germany, Austria, Slovakia, Hungary, Serbia, Bulgaria and Romania to the river's mouth at the Danube Delta. It then continues southwards to end in Constanța, on the Black Sea.
EuroVelo 7 – Sun Route
EV7 runs from the North Cape to Malta. It goes through Norway, Finland, Sweden, Denmark, Germany, Czechia, Austria, Italy, and Malta.
EuroVelo 8 – Mediterranean Route
EV8 follows the European coastline of the Mediterranean sea from Cádiz, Spain to Athens, Greece, going through Spain, France, italy, Slovenia, Croatia, Bosnia and Herzegovina, Montenegro, Albania, Greece, Turkey, and Cyprus.
EuroVelo 9 – Amber Route
EV9 (in Poland, also labeled as R9) stretches from the Baltic Sea to the Adriatic Sea. It is so named after the precious stone amber collected in the Baltic, which was taken by routes such as this to the Mediterranean. One of the shortest of the EuroVelo routes, EV9 still manages to cut across Europe from north to south, from Poland to Croatia, and in doing so passes through the Czech Republic, Austria and Slovenia en route.
EuroVelo 10 – Baltic Route
EV10 runs around Baltic Sea. Some of its parts are mapped on OpenStreetMap project Relation: EuroVelo 10 - Baltic Sea Cycle Route - part Sweden (63584). On the state of the route there is an OpenStreetMap wiki page
EuroVelo 11 – East Europe Route
EV11 connects (theoretically) Norway's North Cape with Athens.
EuroVelo 12 – North Sea Route
EV12 was the first European route, opened in June 2001, route through England, Scotland, Norway, Sweden, Denmark, Germany, the Netherlands and Belgium. It features in the Guinness Book of Records as the longest unbroken signposted cycling route. It was funded in part by the European Union's Interreg initiative.
EuroVelo 13 – Iron Curtain Trail
EuV13 follows the old Iron Curtain, the divided borders of Europe during the Cold War. The ICT runs from Kirkenes, Norway on the Barents Sea, along the Finno-Russian border through to the Baltic Sea, then hugs the length of the Baltic coast to Lübeck in Germany. It then follows the old border between West Germany and the former East Germany, the current borders between the Czech Republic and both Germany then Austria, the Austrian-Slovak and Austrian-Hungarian borders before following the borders of Romania, the former Yugoslavia, Bulgaria and North Macedonia. It finishes at Rezovo in Bulgaria on the Black Sea after following the border with Greece and Turkey.
EuroVelo 15 – The Rhine Cycle Route
EV15, with an overall length of about passes through four countries from the headwaters of the Rhine in Andermatt in the Swiss Alps to the estuary in Rotterdam in the Netherlands, via France and Germany.
EuroVelo 17 – Rhone Cycle Route
EV17 has an overall length of about . It starts in Andermatt and runs along each side of Lake Geneva before crossing into France. Passing through Lyon and Avignon, it forks into sections which end in Montpellier and Marseille.
EuroVelo 19 – Meuse Cycle Route
EV19, with an overall length of about , is the newest and the shortest EuroVelo route. It follows one of the most significant rivers in Europe, from the source of the Meuse on the Langres plateau in France, heading north into Belgium and on to the river mouth at Hook of Holland, with the route ending in the Dutch port city of Rotterdam.
Requirements
The ECF has written a route development manual for those working on developing EuroVelo routes. According to the guidelines, all EuroVelo routes should fulfill the following criteria:
They must be based on existing or planned national or regional routes of the involved countries.
At least two countries must be involved.
Route length must be at least .
Steep sections should be avoided wherever possible and for very steep sections (if unavoidable) alternative transport options (i.e. public transport or alternative routes) should be provided.
Easy to communicate - internationally recognisable identity and name (marketing potential).
Implementation plans in place (project plan, business plan, partners).
Signing in accordance with the regulations of the respective nations and/or regions, continuous and in both directions.
Signage supplemented by EuroVelo route information panels, in accordance with the recommendations of UNECE and the ECF's Signing of EuroVelo cycle routes manual.
Route infrastructure
In 2011 the share of route infrastructure components in the EuroVelo network was as follows:
Bicycle path/lane: 14%
Traffic-free asphalted road: 8%
Traffic-free non-asphalted road: 6%
Public low-traffic, asphalted road: 56%
Public non-asphalted road: 3%
Public high-traffic, asphalted road: 14%
| Technology | Ground transportation networks | null |
758445 | https://en.wikipedia.org/wiki/Baler | Baler | A baler or hay baler is a piece of farm machinery used to compress a cut and raked crop (such as hay, cotton, flax straw, salt marsh hay, or silage) into compact bales that are easy to handle, transport, and store. Often, bales are configured to dry and preserve some intrinsic (e.g. the nutritional) value of the plants bundled. Different types of balers are commonly used, each producing a different type of balerectangular or cylindrical, of various sizes, bound with twine, strapping, netting, or wire.
Industrial balers are also used in material recycling facilities, primarily for baling metal, plastic, or paper for transport.
History
Before the 19th century, hay was cut by hand and most typically stored in haystacks using hay forks to rake and gather the scythed grasses into optimally sized heapsneither too large, promoting conditions favorable for spontaneous combustion, nor too small, which would mean much of the pile is susceptible to rotting. These haystacks lifted most of the plant fibers up off the ground, letting air in and water drain out, so the grasses could dry and cure, to retain nutrition for livestock feed at a later time.
In the 1860s, mechanical cutting devices were developed; from these came modern devices including mechanical mowers and balers. In 1872, a reaper that used a knotter device to bundle and bind hay was invented by Charles Withington; this was commercialized in 1874 by Cyrus McCormick. In 1936, Innes invented an automatic baler that tied bales with twine using Appleby-type knotters from a John Deere grain binder; in 1938, Edwin Nolt filed a patent for an improved version that was more reliable.
The first round baler was probably invented in the late 19th century and one was shown in Paris by Pilter (as illustrated by Michael Williams in Steam Power in Agriculture: Blandford, 1977). This was a portable machine designed for use with threshing machines.
Round baler
The most common type of baler in industrialized countries today is the round baler. It produces cylinder-shaped "round" or "rolled" bales. The design has a "thatched roof" effect that withstands weather. Grass is rolled up inside the baler using rubberized belts, fixed rollers, or a combination of the two. When the bale reaches a predetermined size, either netting or twine is wrapped around it to hold its shape. The back of the baler swings open, and the bale is discharged. The bales are complete at this stage, but they may also be wrapped in plastic sheeting by a bale wrapper, either to keep hay dry when stored outside or convert damp grass into silage. Variable-chamber large round balers typically produce bales from in diameter and up to in width. The bales can weigh anywhere from , depending upon size, material, and moisture content. Common modern small round balers (also called "mini round balers") produce bales in diameter and in width, generally weighing from .
Originally conceived by Ummo Luebben circa 1910, the first round baler did not see production until 1947 when Allis-Chalmers introduced the Roto-Baler. Marketed for the water-shedding and light weight properties of its hay bales, AC had sold nearly 70,000 units by the end of production in 1960.
The next major innovation began in 1965 when a graduate student at Iowa State University, Virgil Haverdink, sought out Wesley F. Buchele, a professor of Agricultural Engineering, seeking a research topic for a master thesis. Over the next year, Buchele and Haverdink developed a new design for a large round baler, completed and tested in 1966, and thereafter dubbed the Buchele–Haverdink large round baler. The large round bales were about in diameter, long, and they weighed about after they driedabout 80 kg/m3 (5 lb/ft3). The design was promoted as a "Whale of a Bale" and Iowa State University now explains the innovative design as follows:
In the summer of 1969, the Australian Econ Fodder Roller baler came out, a design that made a ground-rolled bale. In September of that same year, The Hawkbilt Company of Vinton, Iowa, contacted Dr. Buchele about his design, then fabricated a large ground-rolling round baler which baled hay that had been laid out in a windrow, and began manufacturing large round balers in 1970.
In 1972, Gary Vermeer of Pella, Iowa, designed and fabricated a round baler after the design of the A-C Roto-Baler, and the Vermeer Company began selling its model 605the first modern round baler. The Vermeer design used belts to compact hay into a cylindrical shape as is seen today.
In the early 1980s, collaboration between Walterscheid and Vermeer produced the first effective uses of CV joints in balers, and later in other farm machinery. Due to the heavy torque required for such equipment, double Cardan joints are primarily used. Former Walterscheid engineer Martin Brown is credited with "inventing" this use for universal joints.
By 1975, fifteen American and Canadian companies were manufacturing large round balers.
Transport, handling, and feeding
Short-haul transport and on-field handling
Due to the ability for round bales to roll away on a slope, they require specific treatment for safe transport and handling. Small round bales can typically be moved by hand or with lower-powered equipment. Due to their size and their weight, which can be a ton or more, large round bales require special transport and moving equipment.
The most important tool for large round bale handling is the bale spear or spike, which is usually mounted on the back of a tractor or the front of a skid-steer. It is inserted into the approximate center of the round bale, then lifted and the bale is hauled away. Once at the destination, the bale is set down, and the spear pulled out. Careful placement of the spear in the center is needed or the bale can spin around and touch the ground while in transport, causing a loss of control. When used for wrapped bales that are to be stored further, the spear makes a hole in the wrapping that must be sealed with plastic tape to maintain a hermetic seal.
Alternatively, a grapple fork may be used to lift and transport large round bales. The grapple fork is a hydraulically driven implement attached to the end of a tractor's bucket loader. When the hydraulic cylinder is extended, the fork clamps downward toward the bucket, much like a closing hand. To move a large round bale, the tractor approaches the bale from the side and places the bucket underneath the bale. The fork is then clamped down across the top of the bale, and the bucket is lifted with the bale in tow. Grab hooks installed on the bucket of a tractor are another tool used to handle round bales, and can be done by a farmer with welding skills by welding two hooks and a heavy chain to the outside top of a tractor front loader bucket.
Long-haul transport
The rounded surface of round bales poses a challenge for long-haul, flat-bed transport, as they could roll off of the flat surface if not properly supported. This is particularly the case with large round bales; their size makes them difficult to flip, so it may not be feasible to flip many of them onto the flat surface for transport and then re-position them on the round surface at the destination. One option that works with both large and small round bales is to equip the flat-bed trailer with guard rails at either end, which prevent bales from rolling either forward or backward. Another solution is the saddle wagon, which has closely spaced rounded saddles or support posts in which round bales sit. The tall sides of each saddle prevent the bales from rolling around while on the wagon, as the bale settles down in between posts. On 3 September 2010, on the A381 in Halwell near Totnes, Devon, England, an early member of British rock group ELO Mike Edwards was killed when his van was crushed by a large round bale. The cellist, 62, died instantly when the bale fell from a tractor on nearby farmland before rolling onto the road and crushing his van.
Feeding
A large round bale can be directly used for feeding animals by placing it in a feeding area, tipping it over, removing the bale wrap, and placing a protective ring (a ring feeder) around the outside so that animals will not walk on hay that has been peeled off the outer perimeter of the bale. The round baler's rotational forming and compaction process also enables both large and small round bales to be fed out by unrolling the bale, leaving a continuous flat strip in the field or behind a feeding barrier.
Silage or haylage
Silage, a fermented animal feed, was introduced in the late 1800s, and can also be stored in a silage or haylage bale, which is a high-moisture bale wrapped in plastic film. These are baled much wetter than hay bales, and are usually smaller than hay bales because the greater moisture content makes them heavier and harder to handle. These bales begin to ferment almost immediately, and the metal bale spear stabbed into the core becomes very warm to the touch from the fermentation process.
Silage or haylage bales may be wrapped by placing them on a rotating bale spear mounted on the rear of a tractor. As the bale spins, a layer of plastic cling film is applied to the exterior of the bale. This roll of plastic is mounted in a sliding shuttle on a steel arm and can move parallel to the bale axis, so the operator does not need to hold up the heavy roll of plastic. The plastic layer extends over the ends of the bale to form a ring of plastic approximately wide on the ends, with hay exposed in the center.
To stretch the cling-wrap plastic tightly over the bale, the tension is actively adjusted with a knob on the end of the roll, which squeezes the ends of the roll in the shuttle. In the example wrapping video, the operator is attempting to use high tension to get a flat, smooth seal on the right end. However, the tension increases too much and the plastic tears off. The operator recovers by quickly loosening the tension and allows the plastic to feed out halfway around the bale before reapplying the tension to the sheeting.
These bales are placed in a long continuous row, with each wrapped bale pressed firmly against all the other bales in the row before being set down onto the ground. The plastic wrap on the ends of each bale sticks together to seal out air and moisture, protecting the silage from the elements. The end-bales are hand-sealed with strips of cling plastic across the opening.
The airtight seal between each bale permits the row of round bales to ferment as if they were in a silo bag, but they are easier to handle than a silo bag, as they are more robust and compact. The plastic usage is relatively high, and there is no way to reuse the silage-contaminated plastic sheeting, although it can be recycled or used as a fuel source via incineration. The wrapping cost is approximately US$5 per bale.
An alternative form of wrapping uses the same type of bale placed on a bale wrapper, consisting of pair of rollers on a turntable mounted on the three-point linkage of a tractor. It is then spun about two axes while being wrapped in several layers of cling-wrap plastic film. This covers the ends and sides of the bale in one operation, thus sealing it separately from other bales. The bales are then moved or stacked using a special pincer attachment on the front loader of a tractor, which does not damage the film seal. They can also be moved using a standard bale spike, but this punctures the airtight seal, and the hole in the film must be repaired after each move.
Plastic-wrapped bales must be unwrapped before being fed to livestock to prevent accidental ingestion of the plastic. Like round hay bales, silage bales are usually fed using a ring feeder.
Large square baler
In 1978, Hesston introduced the first "large square baler", capable of compacting hay into more easily transported large square bales that could be stacked and tarped in the field (to protect them from rain) or loaded on trucks or containers for trucking or export. Depending upon the baler, these bales can weigh from for a or bale (versus for a round bale). As the pickup revolves just above the ground surface, the tines pick up and feed the hay into the flake forming chamber, where a "flake" of hay is formed before being pushed up into the path of the plunger, which then compresses it with great force (, depending on model) against the existing bale in the chamber. Once the desired length is achieved, the knotter arm is mechanically tripped to begin the knotting cycle in which several knotters (4–6 is common) tie the 4–6 strings that maintain the bale's shape.
In the prairies of Canada, the large rectangular balers are also called "prairie raptors".
Large square bale handling and transport
Square bales are easier to transport than round bales, since there is little risk of the bale rolling off the back of a flatbed trailer. The rectangular shape also saves space and allows a complete solid slab of hay to be stacked for transport and storage. Most balers allow adjustment of length and it is common to produce bales of twice the width, allowing stacks with brick-like alternating groups overlapping the row below at right angles, creating a strong structure.
They are well-suited for large-scale livestock feedlot or dairy operations, where many tons of feed are rationed every hour. Most often, they are baled small enough that one person can carry or toss them where needed.
Due to the huge rectangular shape, large spear forks, or squeeze grips, are mounted to heavy lifting machinery, such as large forklifts, tractors equipped with front-end loaders, telehandlers, hay squeezes or wheel loaders to lift these bales.
Small square baler
The original type of baler produces small square bales. These bales are rectangular-shaped "square" bales. This was once the most prevalent form of baler but is less common today. It is primarily used on small acreages where large equipment is impractical and hay production for small operations, particularly horse owners who may not have access to the specialized feeding machinery used for larger bales. Each bale is about . The bales are usually wrapped with two, sometimes three, or more strands of knotted twine. The bales are light enough for one person to handle, about , depending upon the crop and pressure applied (can be 100 lbs for a 16"x18" 2-string bale and even more for a 3-string bale). Many balers have adjustable bale chamber pressure and bale length, so shorter, less-dense bales can be produced for easy handling.
To form the bale, the material to be baled (which is often hay or straw) in the windrow is lifted by tines in the baler's reel. This material is then packed into the bale chamber, which runs the length of one side of the baler (usually the left-hand side when viewed from the rear) in offset balers. Balers like Hesston models use an in-line system where the hay goes straight through from the pickup to the flake chamber to the plunger and bale-forming chamber. A combination plunger and knife move back and forth in the front of this chamber, with the knife closing the door into the bale chamber as it moves backward. The plunger and knife are attached to a heavy asymmetrical flywheel to provide extra force as they pack the bales. A measuring device, typically a spiked wheel that is turned by the emerging balesmeasures the amount of material that is being compressed and, at the appropriate length, it triggers the knotters that wrap the twine around the bale and tie it off. As the next bale is formed, the tied one is driven out of the rear of the baling chamber, where it can either drop to the ground or be sent to a wagon or accumulator towed behind the baler.
When a wagon is used, the bale may be lifted by hand from the chamber by a worker on the wagon who stacks the bales on the wagon, or the bale may be propelled into the wagon by a mechanism on the baler, commonly either a "thrower" (parallel high-speed drive belts which throw the bale into the wagon) or a "kicker" (mechanical arm which throws the bale into the wagon). In the case of a thrower or kicker, the wagon has high walls on the left, right, and back sides and a short wall on the front side to contain the randomly piled bales. This process continues as long as the material is in the bale chamber and there is twine to tie the bales.
This form of bale is not used much in large-scale commercial agriculture because the efficiency and speed of large bales are higher. However, it has some popularity in small-scale, low-mechanization agriculture and horse-keeping. Besides using simpler machinery and being easy to handle, these small bales can also be used for insulation and building materials in straw-bale construction. Convenience is also a significant factor in farmers deciding to continue putting up square bales, as they make feeding and bedding in confined areas (stables, barns, etc.) much more manageable and thus command a higher market value per ton.
The automatic baler for small square bales took on most of its present form in 1938, with the first baler sold as Arthur S. Young's Automaton Baler. It was manufactured in small numbers until New Holland Ag acquired it.
In Europe, as early as 1939, both Claas of Germany and Rousseau SA of France had automatic twine-tying pick-up balers. Most of these produced low-density bales, however. The first successful pick-up balers were made by the Ann Arbor Company in 1929. Ann Arbor was acquired by the Oliver Farm Equipment Company in 1943. Despite their head start on the rest of the field, no Ann Arbor balers carried automatic knotters or twisters and Oliver did not produce its own automatic tying baler until 1949.
Small square bale handling and transport
In the 1940s most farmers would bale hay in the field with a small tractor with 20 or less horsepower, and the tied bales would be dropped onto the ground as the baler moved through the field. Another team of workers with horses and a flatbed wagon would come by and use a sharp metal hook to grab the bale and throw it up onto the wagon while an assistant stacked the bales, for transport to the barn.
A later time-saving innovation was to tow the flatbed wagon directly behind the baler. The bale would be pushed up a ramp to a waiting attendant on the wagon. The attendant would hook the bale off the ramp and stack it on the wagon while waiting for the next bale to be produced.
Eventually, as balers evolved, the bale thrower was developed, eliminating needing someone to stand on the wagon and pick up the finished bales. The first thrower mechanism used two fast-moving friction belts to grab finished bales and throw them at an angle up in the air onto the bale wagon. The bale wagon was modified from a flatbed into a three-sided skeleton frame open at the front to act as a catcher's net for the thrown bales.
As tractor horsepower increased, the thrower-baler's next innovation was the hydraulic tossing baler. This employs a flat pan behind the bale knotter. As bales advance out the back of the baler, they are pushed onto the pan one at a time. When the bale has moved entirely onto the pan, the pan suddenly pops up, pushed by a large hydraulic cylinder, and tosses the bale up into the wagon like a catapult.
The pan-thrower method puts much less stress on the bales than the belt-thrower. The friction belts of the belt-thrower stress the twine and knots as they grip the bale and occasionally cause bales to break apart in the thrower or when the bales land in the wagon.
Mechanical small square bale handling
There are several ways to handle small bales automatically and eliminate hand labor almost completely. Bale stackers, bale bundlers, bale accumulators, and bale sledges are the different categories of these machines.
Bale Stackers: Bales may be picked up from the field and stacked using a self-powered machine called a bale stacker, bale wagon or harobed. There are several designs and sizes made by New Holland. One type picks up square bales, which are dropped by the baler with the strings facing sideways. The stacker will drive up to each bale, pick it up, and set it on a three-bale-wide table (the strings are now facing upwards). Once three bales are on the table, the table lifts up and back, causing the three bales to face strings to the side again; this happens three more times until there are 16 bales on the main table. This table will lift like the smaller one, and the bales will be up against a vertical table. The machine will hold 160 bales (ten tiers); usually, there will be cross-tiers near the center to keep the stack from swaying or collapsing if any weight is applied to the top of the stack. The full load will be transported to a barn; the whole rear of the stacker will tilt upwards until it is vertical. There will be two pushers that will extend through the machine and hold the bottom of the stack from being pulled out from the stacker while it is driven out of the barn.
Bale Bundlers: Bales may be picked up from the field or collected directly from the small square baler. Each bale is ingested into a compression chamber and indexed until either two or three are ready for compression. After 7 compression cycles are completed, making a cube of 14 or 21 bales they are tied with twine or banded into a bundle and ejected onto the ground. These bundles are then handled with spears, grabs, or pallet forks. They are of ideal dimensions for filling van trailers.
Bale Accumulators: Typically these are attached directly to the baler and arrange the small square bales into groups to be retrieved with a "bale grabber" or "bale grab" mounted on a loader. There are a number of different methods employed by these machines to arrange the bales into groups. One method is to allow up to three bales to be pushed in line onto a tray that is then emptied sideways by a hydraulically driven push bar onto a platform. After four or five such pushes, a group is made of 8, 10 12, or 15 bales and the platform is emptied onto the ground.
Another, more efficient method that doesn't use hydraulics, is to use a system of levers and gates to guide bales into channels. These can have four, five, or six channels and accommodate two or three bales per channel. This method can make groups of 4, 8, 10, 12, 15 or 18 bales. In this method, the last bale of the group triggers a rear gate open, and the bales are deposited on the ground.
These groups can be bound with twine for stack stability or not and be stacked on wagons or trailers for transport to storage. These groups are ideal for storage in buildings accessible to equipment. This is also the ideal way to automate bales and also allow them to cure properly.
Bale Sledge: In Britain (if small square bales are still to be used), they are usually collected as they fall out of the baler in a bale sledge dragged behind the baler. This has four channels, controlled by automatic mechanical balances, catches, and springs, which sort each bale into its place in a square eight. When the sledge is full, a catch is tripped automatically, and a door at the rear opens to leave the eight lying neatly together on the ground. These may be picked up individually and loaded by hand, or they may be picked up all eight together by a bale grab on a tractor, a special front loader consisting of many hydraulically powered downward-pointing curved spikes. The square eight will then be stacked, either on a trailer for transport or in a roughly cubic field stack eight or ten layers high. This cube may then be transported by a large machine attached to the three-point hitch behind a tractor, which clamps the sides of the cube and lifts it bodily.
Square/wire bale history
Hay presses, wire balers
Prior to 1937, the hay press was the common name of the stationary baling implement, powered with a tractor or stationary engine using a belt on a belt pulley, with the hay being brought to the baler and fed in by hand. Later, balers were made mobile, with a 'pickup' to gather up the hay and feed it into the chamber. These often used air-cooled gasoline engines mounted on the baler for power. The biggest change to this type of baler since 1940 is being powered by the tractor through its power take-off (PTO), instead of by a built-in internal combustion engine.
In present-day production, small square balers can be ordered with twine knotters or wire tie knotters.
Not all stationary wire-tying balers used two wires. It was not uncommon for the larger bale sizeusually machines to use 'boards' that had three slots for wires and hence tied three wires per bale. Most North American manufacturers produced these machines as either regular models or as size options. 'Small square' three wire tying pick-up balers were available from the early 1930s, principally from J. I. Case & Co. and Ann Arbor. These machines were hand-tying and hand-threading machines.
Storage methods
Before electrification occurred in rural parts of the United States in the 1940s, some small dairy farms would have tractors but not electric power. Often just one neighbor who could afford a tractor would do all the baling for surrounding farmers still using horses.
To get the bales up into the hayloft, a pulley system ran on a track along the peak of the barn's hayloft. This track also stuck a few feet out the end of the loft, with a large access door under the track. On the bottom of the pulley system was a bale spear, which was pointed on the end and had retractable retention spikes.
A flatbed wagon would pull up next to the barn underneath the end of the track, the spear lowered down to the wagon and speared into a single bale. The pulley rope would be used to manually lift the bale up until it could enter the mow through the door, then moved along the track into the barn and finally released for manual stacking in tight rows across the floor of the loft. As the stack filled the loft, the bales would be lifted higher and higher with the pulleys until the hay was stacked all the way up to the peak.
When electricity arrived, the bale spear, pulley, and track system were replaced by long motorized bale conveyors known as hay elevators. A typical elevator is an open skeletal frame, with a chain that has dull spikes every few feet along the chain to grab bales and drag them along. One elevator replaced the spear track and ran the entire length of the peak of the barn. A second elevator was either installed at a 30-degree slope on the side of the barn to lift bales up to the peak elevator or used dual front-back chains surrounding the bale to lift bales straight up the side of the barn to the peak elevator.
A bale wagon pulled up next to the lifting elevator, and a farm worker placed bales one at a time onto the angled track. Once bales arrived at the peak elevator, adjustable tipping gates along the length of the peak elevator were opened by pulling a cable from the floor of the hayloft, so that bales tipped off the elevator and dropped down to the floor in different areas of the loft. This permitted a single elevator to transport hay to one part of a loft and straw to another part.
This complete hay elevator lifting, transport, and dropping system reduced bale storage labor to a single person, who simply pulls up with a wagon, turns on the elevators, and starts placing bales on it, occasionally checking to make sure that bales are falling in the right locations in the loft.
The neat stacking of bales in the loft is often sacrificed for the speed of just letting them fall and roll down the growing pile in the loft, and changing the elevator gates to fill in open areas around the loose pile. But if desired, the loose bale pile dropped by the elevator could be rearranged into orderly rows between wagon loads.
Usage once in the barn
The process of retrieving bales from a hayloft has stayed relatively unchanged from the beginning of baling. Typically workers were sent up into the loft, to climb up onto the bale stack, pull bales off the stack, and throw or roll them down the stack to the open floor of the loft. Once the bale is down on the floor, workers climb down the stack, open a cover over a bale chute in the floor of the loft, and push the bales down the chute to the livestock area of the barn.
Most barns were equipped with several chutes along the sides and in the center of the loft floor. This permitted bales to be dropped into the area where they were to be used. Hay bales would be dropped through side chutes, to be broken up and fed to the cattle. Straw bales would be dropped down the center chute, to be distributed as bedding in the livestock standing/resting areas.
Traditionally multiple bales were dropped down to the livestock floor and the twine was removed by hand. After drying and being stored under tons of pressure in the haystack, most bales are tightly compacted and need to be torn apart and fluffed up for use.
One recent method of speeding up all this manual bale handling is the bale shredder, which is a large vertical drum with rotary cutting/ripping teeth at the base of the drum. The shredder is placed under the chute and several bales are dropped in. A worker then pushes the shredder along the barn aisle as it rips up a bale and spews it out in a continuous fluffy stream of material.
Industrial balers
Industrial balers are typically used to compact similar types of waste, such as office paper, cardboard, plastic, foil, and cans, for sale to recycling companies. These balers are made of steel with a hydraulic ram to compress the material loaded. Some balers are simple and labor-intensive but are suitable for smaller volumes. Other balers are very complex and automated and are used where large quantities of waste are handled.
Used in recycling facilities, balers are a packaging step that allows for the aforementioned commodities to be broken down into dense cubes of one type of material at a time. There are different balers used depending on the material type. After a specific material is crushed down into a dense cube, it is tied to a bale by a thick wire and then pushed out of the machine. This process allows for easy transport of all materials involved.
Two-ram baler: A two-ram baler is a baling machine that contains two cylinders and is able to bundle and package most commodities except for cardboard and clear film. This baler is known for its durability and is able to take in more bulky material.
Single-ram baler: A single-ram baler is a baling machine that contains one cylinder. Because this baler is relatively smaller than the two-ram baler, it is best for small and medium commodities.
Closed door baler: This baler bales clear plastic film.
American baler: This baler bales corrugated materials.
| Technology | Farm and garden machinery | null |
758833 | https://en.wikipedia.org/wiki/Past | Past | The past is the set of all events that occurred before a given point in time. The past is contrasted with and defined by the present and the future. The concept of the past is derived from the linear fashion in which human observers experience time, and is accessed through memory and recollection. In addition, human beings have recorded the past since the advent of written language.
In English, the word past was one of the many variant forms and spellings of passed, the past participle of the Middle English verb passen (whence Modern English pass), among ypassed, ypassyd, i-passed, passyd, passid, pass'd, paste, etc. It developed into an adjective and preposition in the 14th century, and a noun (as in the past or a past, through ellipsis with the adjective past) in the 15th century.
Grammar
In English grammar, actions are classified according to one of the following twelve verb tenses: past (past, uses of English verb forms, past perfect, or past perfect continuous), present (present, present continuous, present perfect, or present perfect continuous), or future (future, future continuous, future perfect, or future perfect continuous). The past tense refers to actions that have already happened. For example, "she is walking" refers to a girl who is currently walking (present tense), while "she walked" refers to a girl who was walking before now (past tense).
The past continuous tense refers to actions that continued for a period of time, as in the sentence "she was walking," which describes an action that was still happening in a prior window of time to which a speaker is presently referring. The past perfect tense is used to describe actions that were already completed by a specific point in the past. For example, "she had walked" describes an action that took place in the past and was also completed in the past.
The past perfects continuous tense refers to an action that was happening up until a particular point in the past but was completed. It is different from the past perfect tense because the emphasis of past perfect continuous verbs is not on the action having been completed by the present moment, but rather on its having taken place actively over a time period before another moment in the past. The verb tense used in the sentence "She had been walking in the park regularly before I met her" is past perfect continuous because it describes an action ("walking") that was actively happening before a time when something else in the past was happening (when "I met her").
Depending on its usage in a sentence, "past" can be described using a variety of terms. Synonyms for "past" as an adjective include, "former," "bygone," "earlier," "preceding," and "previous." Synonyms for "past" as a noun include, "history, "background," "life story," and "biography." Synonyms of "past" as a preposition include, "in front of," "beyond," "by," and "in excess of."
Other uses
The word "past" can also be used to describe the offices of those who have previously served in an organization, group, or event such as, "past president," or, "past champions." "Past" can also refer to something or someone being at or in a position that is further than a particular point. For instance, in the sentence, "I live on Fielding Road, just past the train station," the word "past" is used to describe a location (the speaker's residence) beyond a certain point (the train station). Alternatively, the sentence, "He ran past us at full speed," utilizes the concept of the past to describe the position of someone ("He") that is further than the speaker.
The "past" is also used to define a time that is a certain number of minute before or after a particular hour, as in "We left the party at half-past twelve." People also use "past" to refer to being beyond a particular biological age or phase of being, as in, "The boy was past the age of needing a babysitter," or, "I'm past caring about that problem." The "past" is commonly used to refer to history, either generally or with regards to specific time periods or events, as in, "Past monarchs had absolute power to determine the law in contrast to many European Kings and Queens of today."
Nineteenth-century British author Charles Dickens created one of the best-known fictional personifications of the "past" in his short book, "A Christmas Carol." In the story, the Ghost of Christmas Past is an apparition that shows the main character, a cold-hearted and tight-fisted man named Ebenezer Scrooge, vignettes from his childhood and early adult life to teach him that joy does not necessarily come from wealth.
Fields of study
The past is the object of study within such fields as time, life, history, nostalgia, archaeology, archaeoastronomy, chronology, geology, historical geology, historical linguistics, ontology, paleontology, paleobotany, paleoethnobotany, palaeogeography, paleoclimatology, etymology and physical cosmology.
| Technology | Timekeeping | null |
759008 | https://en.wikipedia.org/wiki/Eye%20color | Eye color | Eye color is a polygenic phenotypic trait determined by two factors: the pigmentation of the eye's iris and the frequency-dependence of the scattering of light by the turbid medium in the stroma of the iris.
In humans, the pigmentation of the iris varies from light brown to black, depending on the concentration of melanin in the iris pigment epithelium (located on the back of the iris), the melanin content within the iris stroma (located at the front of the iris), and the cellular density of the stroma. The appearance of blue, green, and hazel eyes results from the Tyndall scattering of light in the stroma, a phenomenon similar to Rayleigh scattering which accounts for the blue sky. Neither blue nor green pigments are present in the human iris or vitreous humour. This is an example of structural color, which depends on the lighting conditions, especially for lighter-colored eyes.
The brightly colored eyes of many bird species result from the presence of other pigments, such as pteridines, purines, and carotenoids. Humans and other animals have many phenotypic variations in eye color.
The genetics and inheritance of eye color in humans is complicated. , as many as 16 genes have been associated with eye color inheritance. Some of the eye-color genes include OCA2 and HERC2. The earlier belief that blue eye color is a recessive trait has been shown to be incorrect, and the genetics of eye color are so complex that almost any parent-child combination of eye colors can occur.
Genetic determination
Eye color is an inherited trait determined by multiple genes. These genes are sought by studying small changes in the genes themselves and in neighboring genes, called single-nucleotide polymorphisms or SNPs. The total number of genes that contribute to eye color is unknown, but there are a few likely candidates. A study in Rotterdam (2009) found that it was possible to predict eye color with more than 90% accuracy for brown and blue using just six SNPs.
In humans, eye color is a highly sexually dimorphic trait. Several studies have shown that men are more likely to have blue eyes than women, while women are more likely to have darker eye colors (green and brown eyes) than men. Sex is therefore a major factor in the expression of eye color genotypes. One study suggested that women's higher levels of the sex hormone estrogen may explain why women tend to have darker eyes than men.
People of European descent show the greatest variety in eye color of any population worldwide. Recent advances in ancient DNA technology have revealed some of the history of eye color in Europe. Through the analysis of ancient DNA, a 2020 study published in Experimental Dermatology suggested that the common gene for blue eye color likely originated in the Near East and arrived in Europe around 42,000 years ago, after the exodus out of Africa.
There is evidence that as many as 16 different genes could be responsible for eye color in humans; however, the main two genes associated with eye color variation are OCA2 and HERC2, and both are localized in chromosome 15.
The gene OCA2 (), when in a variant form, causes the pink eye color and hypopigmentation common in human albinism. (The name of the gene is derived from the disorder it causes, oculocutaneous albinism type II.) Different SNPs within OCA2 are strongly associated with blue and green eyes as well as variations in freckling, mole counts, hair and skin tone. The polymorphisms may be in an OCA2 regulatory sequence, where they may influence the expression of the gene product, which in turn affects pigmentation. A specific mutation within the HERC2 gene, a gene that regulates OCA2 expression, is partly responsible for blue eyes. Other genes implicated in eye color variation are SLC24A4 and TYR. A 2010 study of eye color variation in hue and saturation values using high-resolution digital full-eye photographs found three new loci for a total of ten genes, allowing the explanation of about 50% of eye color variation.
Blue eyes with a brown spot, green eyes, and gray eyes are caused by an entirely different part of the genome.
Changes in eye color
A 1997 study of White Americans found that eye color may be subject to change in infancy, and from adolescence to adulthood. 17% of children experienced a change of eye color by adulthood. Of those children, 50% of developed lighter eyes as they got older. The other 50% developed darker eyes.
Generally, children with hazel and light brown eyes tended to experience a lightening of their eye color by adulthood. Children with green eyes often experienced a darkening of their eye color. It was also found that 11% of the children's mothers experienced an eye color change during the same period, with most developing lighter eyes, relative to their original color at the time of their child's birth.
Eye color range
Brown
Almost all mammals have brown or darkly-pigmented irises. In humans, brown is by far the most common eye color, with approximately 79% of people in the world having it. Brown eyes result from a relatively high concentration of melanin in the stroma of the iris, which causes light of both shorter and longer wavelengths to be absorbed.
In many parts of the world, it is nearly the only iris color present. Brown eyes are common in Europe, East Asia, Southeast Asia, Central Asia, South Asia, West Asia, Oceania, West Africa and the Americas. Light or medium-pigmented brown eyes can also be commonly found in Europe, among the Americas, and parts of Central Asia, West Asia, South Asia, and East Africa. Light brown eyes bordering amber and hazel coloration are more common in Europe, but can also be observed in East Asia, Southeast Asia, North Africa and East Africa.
Amber
Amber eyes are a solid color with a strong yellowish/golden or russet/coppery tint, which may be due to a yellow pigment called lipochrome (also found in green eyes). Amber eyes should not be confused with hazel eyes. Although hazel eyes may contain specks of amber or gold, they usually tend to have many other colors, including green, brown, and orange. Also, hazel eyes may appear to shift in color and consist of flecks and ripples, while amber eyes are of a solid gold hue. Even though amber is similar to gold, some people have russet- or copper-colored amber eyes that are mistaken for hazel, though hazel tends to be duller and contains green with red/gold flecks, as mentioned above. Amber eyes may also contain amounts of very light gold-ish gray.
The eyes of some pigeons contain yellow fluorescing pigments known as pteridines. The bright yellow eyes of the great horned owl are thought to be due to the presence of the pteridine pigment xanthopterin within certain chromatophores (called xanthophores) located in the iris stroma. In humans, yellowish specks or patches are thought to be due to the pigment lipofuscin, also known as lipochrome. Many animals such as canines, domestic cats, owls, eagles, pigeons, and fish have amber eyes, whereas in humans this color occurs less frequently. Amber is the third-rarest natural eye color after green and gray, occurring in 5% of the world's population. People with amber-colored eyes are found in Europe, and in fewer numbers in the Middle East, North Africa, and South America.
Hazel
Hazel eyes are due to a combination of Rayleigh scattering and a moderate amount of melanin in the iris' anterior border layer. Hazel eyes often appear to shift in color from a brown to a green. Although hazel mostly consists of brown and green, the dominant color in the eye can either be brown/gold or green. This is why hazel eyes can be mistaken as amber, and why amber is often counted as hazel in studies, and vice versa. The combination can sometimes produce a multicolored iris, i.e. an eye that is light brown/amber near the pupil and charcoal or dark green on the outer part of the iris (or vice versa) when observed in sunlight.
Definitions of the eye color "hazel" vary: it is sometimes considered to be synonymous with light brown or gold, as in the color of a hazelnut shell.
Around 18% of the US population and 5% of the world population have hazel eyes. 55.2% of Spanish subjects in a series of 221 photographs were judged to have hazel eyes. Hazel eyes are found in Europe, most commonly in the Netherlands and the United Kingdom, and have also been observed to be very common among the Low Saxon-speaking populations of northern Germany.
Green
Green eyes probably result from the interaction of multiple allelic variants of OCA2 and other genes. They may have been present in southern Siberia during the Bronze Age.
Green eyes are most common in Northern, Western, and Central Europe. Around 8–10% of men and 18–21% of women in Iceland and 6% of men and 17% of women in the Netherlands have green eyes. Among European Americans, green eyes are most common among those of recent Celtic and Germanic ancestry, occurring in about 16% of people with those backgrounds.
The green color is caused by the combination of: 1) an amber or light brown pigmentation in the stroma of the iris (which has a low or moderate concentration of melanin), and 2) a blue shade created by the Rayleigh scattering of reflected light. Green eyes contain the yellowish pigment lipochrome.
Blue
There is no intrinsically blue pigmentation either in the iris or in the vitreous body. Rather, blue eyes result from structural color in combination with certain concentrations of non-blue pigments. The iris pigment epithelium is brownish black due to the presence of melanin. Unlike brown eyes, blue eyes have low concentrations of melanin in the stroma of the iris, which lies in front of the dark epithelium. Longer wavelengths of light tend to be absorbed by the dark underlying epithelium, while shorter wavelengths are reflected and undergo Rayleigh scattering in the turbid medium of the stroma. This is the same scattering that accounts for the blue appearance of the sky. The result is a "Tyndall blue" structural color that varies with external lighting conditions.
Blue eyes are a highly sexually dimorphic eye color. Studies from various populations in Europe have shown that men are substantially more likely to have blue eyes than women.
The inheritance pattern followed by blue eyes was previously assumed to be a Mendelian recessive trait, though this has been shown to be incorrect. Eye color inheritance is now recognized as a polygenic trait, meaning that it is controlled by the interactions of several genes.
In 2008, a team of researchers from the University of Copenhagen located a single mutation that causes the phenomenon of blue eyes. The research was published in the Journal of Human Genetics. The same DNA sequence of the OCA2 gene among blue-eyed people suggests they may have a single common ancestor. The researchers hypothesized that the OCA2 mutation responsible for blue eyes arose in an individual who lived in the northwestern part of the Black Sea region in Europe sometime between 6,000 and 10,000 years ago, during the Neolithic period. However, more recent ancient DNA research has identified human remains much older than the Neolithic period which possess the OCA2 mutation for blue eyes. It is now believed that the OCA2 allele responsible for blue eyes dates back to the migration of modern humans out of Africa roughly 50,000 years ago, and entered Europe from western Asia.
Eiberg and colleagues suggested in a study published in Human Genetics that a mutation in the 86th intron of the HERC2 gene, which is hypothesized to interact with the OCA2 gene promoter, reduced expression of OCA2 with subsequent reduction in melanin production.
It has been proposed that blue eyes may have been adaptive to shorter day lengths at higher latitudes, as blue eyes increase intraocular light scattering, which suppresses melatonin release from the pineal gland, perhaps reducing psychological depression (which is linked to the short day length of higher latitudes).
Blue eyes are predominant in northern and eastern Europe, particularly around the Baltic Sea. Blue eyes are also found in Southern Europe, Central Asia, South Asia, North Africa, and West Asia.
Approximately 8% to 10% of the global population have blue eyes. A 2002 study found the prevalence of blue eye color among the white population in the United States to be 33.8% for those born from 1936 through 1951, compared with 57.4% for those born from 1899 through 1905. , one out of every six Americans, or 16.6% of the total US population, has blue eyes, including 22.3% of whites. The incidence of blue eyes continues to decline among American children. Of Slovenes, 56% have blue/green eyes. In a series of 221 photographs of Spanish subjects, 16.3% of the subjects were determined to have blue-gray eyes.
Gray
Like blue eyes, gray eyes have a dark epithelium at the back of the iris and a relatively clear stroma at the front. One possible explanation for the difference in appearance between gray and blue eyes is that gray eyes have larger deposits of collagen in the stroma, so that the light that is reflected from the epithelium undergoes Mie scattering (which is not strongly frequency-dependent) rather than Rayleigh scattering (in which shorter wavelengths of light are scattered more). This would be analogous to the change in the color of the sky, from the blue given by the Rayleigh scattering of sunlight by small gas molecules when the sky is clear, to the gray caused by Mie scattering of large water droplets when the sky is cloudy. Alternatively, it has been suggested that gray and blue eyes might differ in the concentration of melanin at the front of the stroma.
Gray eyes can also be found among the Algerian Shawia people of the Aurès Mountains in Northwest Africa, in the Middle East/West Asia, Central Asia, and South Asia. In the Iliad, the Greek goddess Athene is said to have gray eyes (γλαυκῶπις). Under magnification, gray eyes exhibit small amounts of yellow and brown color in the iris.
Gray is the second-rarest natural eye color after green, with 3% of the world's population having it.
Special cases
Two different colors
As a result of heterochromia iridum, it is also possible to have two different eye colors. This occurs in humans and certain breeds of domesticated animals and affects less than 1 percent of the world's population.
Red and violet
The eyes of people with severe forms of albinism may appear red under certain lighting conditions owing to the extremely low quantities of melanin, allowing the blood vessels to show through. In addition, flash photography can sometimes cause a "red-eye effect", in which the very bright light from a flash reflects off the retina, which is abundantly vascular, causing the pupil to appear red in the photograph.
Although the deep blue eyes of some people such as Elizabeth Taylor can appear purple or violet at certain times, "true" violet-colored eyes occur only due to albinism. Eyes that appear red or violet under certain conditions due to albinism occur in less than 1 percent of the world's population.
Medical implications
The most important role of melanin in the iris is to protect the eyes from the sun's harmful rays. People with lighter eye colors, such as blue or green, have lessened protection from the sun, and so need greater protection from the sun's rays than those with darker eye colors.
Those with lighter iris color have been found to have a higher prevalence of age-related macular degeneration (ARMD) than those with darker iris color; lighter eye color is also associated with an increased risk of ARMD progression. A gray iris may indicate the presence of a uveitis, and an increased risk of uveal melanoma has been found in those with blue, green or gray eyes. However, a study in 2000 suggests that people with dark brown eyes are at increased risk of developing cataracts and therefore should protect their eyes from direct exposure to sunlight.
Wilson's disease
Wilson's disease involves a mutation of the gene coding for the enzyme ATPase 7B, which prevents copper within the liver from entering the Golgi apparatus in cells. Instead, the copper accumulates in the liver and in other tissues, including the iris of the eye. This results in the formation of Kayser–Fleischer rings, which are dark rings that encircle the periphery of the iris.
Coloration of the sclera
Eye color outside of the iris may also be symptomatic of disease. Yellowing of the sclera (the "whites of the eyes") is associated with jaundice, and may be symptomatic of liver diseases such as cirrhosis or hepatitis. A blue coloration of the sclera may also be symptomatic of disease.
Aniridia
Aniridia is a congenital condition characterized by an extremely underdeveloped iris, which appears absent on superficial examination.
Ocular albinism and eye color
Normally, there is a thick layer of melanin on the back of the iris. Even people with the lightest blue eyes, with no melanin on the front of the iris at all, have dark brown coloration on the back of it, to prevent light from scattering around inside the eye. In those with milder forms of albinism, the color of the iris is typically blue but can vary from blue to brown. In severe forms of albinism, there is no pigment on the back of the iris, and light from inside the eye can pass through the iris to the front. In these cases, the only color seen is the red from the hemoglobin of the blood in the capillaries of the iris. Such albinos have pink eyes, as do albino rabbits, mice, or any other animal with a total lack of melanin. Transillumination defects can almost always be observed during an eye examination due to lack of iridial pigmentation. The ocular albino also lacks normal amounts of melanin in the retina as well, which allows more light than normal to reflect off the retina and out of the eye. Because of this, the pupillary reflex is much more pronounced in albino individuals, and this can emphasize the red eye effect in photographs.
Heterochromia
Heterochromia (heterochromia iridum or heterochromia iridis) is an eye condition in which one iris is a different color from the other (complete heterochromia), or where a part of one iris is a different color from the remainder (partial heterochromia or sectoral heterochromia). It is a result of the relative excess or lack of pigment within an iris or part of an iris, which may be inherited or acquired by disease or injury. This uncommon condition usually results due to uneven melanin content. A number of causes are responsible, including genetic, such as chimerism, Horner's syndrome and Waardenburg syndrome.
A chimera can have two different colored eyes just like any two siblings can—because each cell has different eye color genes. A mosaic can have two different colored eyes if the DNA difference happens to be in an eye-color gene.
There are many other possible reasons for having two different-colored eyes. For example, the film actor Lee Van Cleef was born with one blue eye and one green eye, a trait that reportedly was common in his family, suggesting that it was a genetic trait. This anomaly, which film producers thought would be disturbing to film audiences, was "corrected" by having Van Cleef wear brown contact lenses. David Bowie, on the other hand, had the appearance of different eye colors due to an injury that caused one pupil to be permanently dilated.
Another hypothesis about heterochromia is that it can result from a viral infection in utero affecting the development of one eye, possibly through some sort of genetic mutation. Occasionally, heterochromia can be a sign of a serious medical condition.
A common cause in females with heterochromia is X-inactivation, which can result in a number of heterochromatic traits, such as calico cats. Trauma and certain medications, such as some prostaglandin analogues, can also cause increased pigmentation in one eye. On occasion, a difference in eye color is caused by blood staining the iris after injury.
Limbal ring
The limbal ring is also a feature of the iris contributing to eye color, appearing as a darkened, occasionally black region encircling the iris resulting from a manifestation of the optical properties of the corneal limbus. Limbal rings are not present in all individuals, and their thickness and prominence may correlate with health or youthfulness, and contributes to facial attractiveness.
Impact on vision
Although people with lighter eye color are generally more sensitive to light because they have less pigment in the iris to protect them from sunlight, there is little to no evidence that eye color has a direct impact on vision qualities such as visual acuity. However, there is a study that found that dark-eyed people perform better at "reactive-type tasks", which suggests they may have better reaction times. People with light-colored eyes, however, performed better at so-called "self-paced tasks", which include activities like hitting a golf ball or throwing baseballs. In another study, people with darker eyes performed better at hitting racquetballs. There are also other studies that challenge these findings, and more study is needed to verify these results.
Classification of color
Iris color can provide a large amount of information about a person, and a classification of colors may be useful in documenting pathological changes or determining how a person may respond to ocular pharmaceuticals. Classification systems have ranged from a basic light or dark description to detailed gradings employing photographic standards for comparison. Others have attempted to set objective standards of color comparison.
The Martin–Schultz scale, developed from the Martin scale, is one standard color scale used in physical anthropology to establish the eye color of an individual. It was created by the anthropologists Rudolf Martin and Bruno K Schultz in the first half of the 20th century. The scale consists of 20 colors ranging from light blue to dark brown-black, corresponding to natural eye colors caused by the amount of melanin in the iris:
Normal eye colors range from the darkest shades of brown to the lightest tints of blue. To meet the need for standardized classification, at once simple yet detailed enough for research purposes, Seddon et al. developed a graded system based on the predominant iris color and the amount of brown or yellow pigment present. There are three pigment colors that determine, depending on their proportion, the outward appearance of the iris, along with structural color. Green irises, for example, have some yellow and the blue structural color. Brown irises contain more or less melanin. Some eyes have a dark ring around the iris, called a limbal ring.
Eye color in non-human animals is regulated differently. For example, instead of blue as in humans, autosomal recessive eye color in the skink species Corucia zebrata is black, and the autosomal dominant color is yellow-green.
As the perception of color depends on viewing conditions (e.g., the amount and kind of illumination, as well as the hue of the surrounding environment), so does the perception of eye color.
| Biology and health sciences | Visual system | Biology |
759071 | https://en.wikipedia.org/wiki/Noctilucent%20cloud | Noctilucent cloud | Noctilucent clouds (NLCs), or night shining clouds, are tenuous cloud-like phenomena in the upper atmosphere of Earth. When viewed from space, they are called polar mesospheric clouds (PMCs), detectable as a diffuse scattering layer of water ice crystals near the summer polar mesopause. They consist of ice crystals and from the ground are only visible during astronomical twilight. Noctilucent roughly means "night shining" in Latin. They are most often observed during the summer months from latitudes between ±50° and ±70°. Too faint to be seen in daylight, they are visible only when the observer and the lower layers of the atmosphere are in Earth's shadow, but while these very high clouds are still in sunlight. Recent studies suggest that increased atmospheric methane emissions produce additional water vapor through chemical reactions once the methane molecules reach the mesosphere – creating, or reinforcing existing, noctilucent clouds.
General
No confirmed record of their observation exists before 1885, although they may have been observed a few decades earlier by Thomas Romney Robinson in Armagh.
Formation
Noctilucent clouds are composed of tiny crystals of water ice up to 100 nm in diameter and exist at a height of about , higher than any other clouds in Earth's atmosphere. Clouds in the Earth's lower atmosphere form when water collects on particles, but mesospheric clouds may form directly from water vapour in addition to forming on dust particles.
Data from the Aeronomy of Ice in the Mesosphere satellite suggests that noctilucent clouds require water vapour, dust, and very cold temperatures to form. The sources of both the dust and the water vapour in the upper atmosphere are not known with certainty. The dust is believed to come from micrometeors, although particulates from volcanoes and dust from the troposphere are also possibilities. The moisture could be lifted through gaps in the tropopause, as well as forming from the reaction of methane with hydroxyl radicals in the stratosphere.
The exhaust from Space Shuttles, in use between 1981 and 2011, which was almost entirely water vapour after the detachment of the Solid Rocket Booster at a height of about , was found to generate minuscule individual clouds. About half of the vapour was released into the thermosphere, usually at altitudes of . In August 2014, a SpaceX Falcon 9 also caused noctilucent clouds over Orlando, Florida after a launch.
The exhaust can be transported to the Arctic region in little over a day, although the exact mechanism of this very high-speed transfer is unknown. As the water migrates northward, it falls from the thermosphere into the colder mesosphere, which occupies the region of the atmosphere just below. Although this mechanism is the cause of individual noctilucent clouds, it is not thought to be a major contributor to the phenomenon as a whole.
As the mesosphere contains very little moisture, approximately one hundred millionth that of air from the Sahara, and is extremely thin, the ice crystals can form only at temperatures below about . This means that noctilucent clouds form predominantly during summer when, counterintuitively, the mesosphere is coldest as a result of seasonally varying vertical winds, leading to cold summertime conditions in the upper mesosphere (upwelling, resulting in adiabatic cooling) and wintertime heating (downwelling, resulting in adiabatic heating). Therefore, they cannot be observed (even if they are present) inside the Polar circles because the Sun is never low enough under the horizon at this season at these latitudes. Noctilucent clouds form mostly near the polar regions, because the mesosphere is coldest there. Clouds in the southern hemisphere are about higher than those in the northern hemisphere.
Ultraviolet radiation from the Sun breaks water molecules apart, reducing the amount of water available to form noctilucent clouds. The radiation is known to vary cyclically with the solar cycle and satellites have been tracking the decrease in brightness of the clouds with the increase of ultraviolet radiation for the last two solar cycles. It has been found that changes in the clouds follow changes in the intensity of ultraviolet rays by about a year, but the reason for this long lag is not yet known.
Noctilucent clouds are known to exhibit high radar reflectivity, in a frequency range of 50 MHz to 1.3 GHz. This behaviour is not well understood but a possible explanation is that the ice grains become coated with a thin metal film composed of sodium and iron, which makes the cloud far more reflective to radar, although this explanation remains controversial. Sodium and iron atoms are stripped from incoming micrometeors and settle into a layer just above the altitude of noctilucent clouds, and measurements have shown that these elements are severely depleted when the clouds are present. Other experiments have demonstrated that, at the extremely low temperatures of a noctilucent cloud, sodium vapour can rapidly be deposited onto an ice surface.
Discovery and investigation
Noctilucent clouds are first known to have been observed in 1885, two years after the 1883 eruption of Krakatoa. It remains unclear whether their appearance had anything to do with the volcanic eruption or whether their discovery was due to more people observing the spectacular sunsets caused by the volcanic debris in the atmosphere. Studies have shown that noctilucent clouds are not caused solely by volcanic activity, although dust and water vapour could be injected into the upper atmosphere by eruptions and contribute to their formation. Scientists at the time assumed the clouds were another manifestation of volcanic ash, but after the ash had settled out of the atmosphere, the noctilucent clouds persisted. Finally, the theory that the clouds were composed of volcanic dust was disproved by Malzev in 1926. In the years following their discovery, the clouds were studied extensively by Otto Jesse of Germany, who was the first to photograph them, in 1887, and seems to have been the one to coin the term "noctilucent cloud". His notes provide evidence that noctilucent clouds first appeared in 1885. He had been doing detailed observations of the unusual sunsets caused by the Krakatoa eruption the previous year and firmly believed that, if the clouds had been visible then, he would undoubtedly have noticed them. Systematic photographic observations of the clouds were organized in 1887 by Jesse, Foerster, and Stolze and, after that year, continuous observations were carried out at the Berlin Observatory.
In the decades after Otto Jesse's death in 1901, there were few new insights into the nature of noctilucent clouds. Wegener's conjecture, that they were composed of water ice, was later shown to be correct. Study was limited to ground-based observations and scientists had very little knowledge of the mesosphere until the 1960s, when direct rocket measurements began. These showed for the first time that the clouds' occurrence coincided with very low temperatures in the mesosphere.
Noctilucent clouds were first detected from space by an instrument on the OGO-6 satellite in 1972. The OGO-6 observations of a bright scattering layer over the polar caps were identified as poleward extensions of these clouds. A later satellite, the Solar Mesosphere Explorer, mapped the distribution of the clouds between 1981 and 1986 with its ultraviolet spectrometer. The clouds were detected with a lidar in 1995 at Utah State University, even when they were not visible to the naked eye. The first physical confirmation that water ice is indeed the primary component of noctilucent clouds came from the HALOE instrument on the Upper Atmosphere Research Satellite in 2001.
In 2001, the Swedish Odin satellite performed spectral analyses on the clouds, and produced daily global maps that revealed large patterns in their distribution.
The AIM (Aeronomy of Ice in the Mesosphere) satellite was launched on 25 April 2007. It was the first satellite dedicated to studying noctilucent clouds, and made its first observations a month later (25 May). Images taken by the satellite show shapes in the clouds that are similar to shapes in tropospheric clouds, hinting at similarities in their dynamics.
In the previous year, scientists with the Mars Express mission had announced their discovery of carbon dioxide–crystal clouds on Mars that extended to above the planet's surface. These are the highest clouds discovered over the surface of a rocky planet. Like noctilucent clouds on Earth, they can be observed only when the Sun is below the horizon.
Research published in the journal Geophysical Research Letters in June 2009 suggests that noctilucent clouds observed following the Tunguska Event of 1908 are evidence that the impact was caused by a comet.
The United States Naval Research Laboratory (NRL) and the United States Department of Defense Space Test Program (STP) conducted the Charged Aerosol Release Experiment (CARE) on September 19, 2009, using exhaust particles from a Black Brant XII suborbital sounding rocket launched from NASA's Wallops Flight Facility to create an artificial noctilucent cloud. The cloud was to be observed over a period of weeks or months by ground instruments and the Spatial Heterodyne IMager for MEsospheric Radicals (SHIMMER) instrument on the NRL/STP STPSat-1 spacecraft. The rocket's exhaust plume was observed and reported to news organizations in the United States from New Jersey to Massachusetts.
A 2018 experiment briefly created noctilucent clouds over Alaska, allowing ground-based measurements and experiments aimed at verifying computer simulations of the phenomenon. A suborbital NASA rocket was launched on 26 January 2018 by University of Alaska professor Richard Collins. It carried water-filled canisters, which were released at about above the Earth. Since the naturally-occurring clouds only appear in summer, this experiment was conducted in mid-winter to assure that its results would not be mixed with a natural event.
Description from satellites
PMC's have four major types based on physical structure and appearance. Type I veils are very tenuous and lack well-defined structure, somewhat like cirrostratus or poorly defined cirrus. Type II bands are long streaks that often occur in groups arranged roughly parallel to each other. They are usually more widely spaced than the bands or elements seen with cirrocumulus clouds. Type III billows are arrangements of closely spaced, roughly parallel short streaks that mostly resemble cirrus. Type IV whirls are partial or, more rarely, complete rings of cloud with dark centres.
Satellite observations allow the very coldest parts of the polar mesosphere to be observed, all the way to the geographic pole. In the early 1970s, visible airglow photometers first scanned the atmospheric horizon throughout the summer polar mesospause region. This experiment, which flew on the OGO-6 satellite, was the first to trace noctilucent-like cloud layers across the polar cap. The very bright scattering layer was seen in full daylight conditions, and was identified as the poleward extension of noctilucent clouds. In the early 1980s, the layer was observed again from a satellite, the Solar Mesospheric Explorer (SME). On board this satellite was an ultraviolet spectrometer, which mapped the distributions of clouds over the time period 1981 to 1986. The experiment measured the altitude profile of scattering from clouds at two spectral channels (primarily) 265 nm and 296 nm.
Polar mesospheric clouds generally increase in brightness and occurrence frequency with increasing latitude, from about 60° to the highest latitudes observed (85°). So far, no apparent dependence on longitude has been found, nor is there any evidence of a dependence on auroral activity.
On 8 July 2018, NASA launched a giant balloon from Esrange, Sweden which traveled through the stratosphere across the Arctic to Western Nunavut, Canada in five days. The giant balloon was loaded with cameras, which captured six million high-resolution images filling up 120 terabytes of data storage, aiming to study the PMCs which are affected by the atmospheric gravity waves, resulted from air being pushed up by mountain ranges all the way up to the mesosphere. These images would aid in studying turbulence in the atmosphere, and consequently better weather forecasting.
NASA uses AIM satellite to study these noctilucent clouds, which always occur during the summer season near the poles. However, tomographic analyses of AIM satellite indicate that there is a spatial negative correlation between albedo and wave‐induced altitude.
Observation
Noctilucent clouds are generally colourless or pale blue, although occasionally other colours including red and green have been observed. The characteristic blue colour comes from absorption by ozone in the path of the sunlight illuminating the noctilucent cloud. They can appear as featureless bands, but frequently show distinctive patterns such as streaks, wave-like undulations, and whirls. They are considered a "beautiful natural phenomenon". Noctilucent clouds may be confused with cirrus clouds, but appear sharper under magnification. Those caused by rocket exhausts tend to show colours other than silver or blue, because of iridescence caused by the uniform size of the water droplets produced.
Noctilucent clouds may be seen at latitudes of 50° to 65°. They seldom occur at lower latitudes (although there have been sightings as far south as Paris, Utah, Italy, Turkey and Spain), and closer to the poles it does not get dark enough for the clouds to become visible. They occur during summer, from mid-May to mid-August in the northern hemisphere and between mid-November and mid-February in the southern hemisphere. They are very faint and tenuous, and may be observed only in twilight around sunrise and sunset when the clouds of the lower atmosphere are in shadow, but the noctilucent cloud is illuminated by the Sun. They are best seen when the Sun is between 6° and 16° below the horizon. Although noctilucent clouds occur in both hemispheres, they have been observed thousands of times in the northern hemisphere, but fewer than 100 times in the southern. Southern hemisphere noctilucent clouds are fainter and occur less frequently; additionally the southern hemisphere has a lower population and less land area from which to make observations.
These clouds may be studied from the ground, from space, and directly by sounding rocket. Also, some noctilucent clouds are made of smaller crystals, 30 nm or less, which are invisible to observers on the ground because they do not scatter enough light.
Forms
The clouds may show a large variety of different patterns and forms. An identification scheme was developed by Fogle in 1970 that classified five different forms. These classifications have since been modified and subdivided.
Type I veils are very tenuous and lack well-defined structure, somewhat like cirrostratus or poorly defined cirrus.
Type II bands are long streaks that often occur in roughly parallel groups, usually more widely spaced than the bands or elements seen with cirrocumulus clouds.
Type III billows are arrangements of closely spaced, roughly parallel short streaks that mostly resemble cirrus.
Type IV whirls are partial or, more rarely, complete rings of cloud with dark centres.
| Physical sciences | Clouds | Earth science |
759264 | https://en.wikipedia.org/wiki/Amount%20of%20substance | Amount of substance | In chemistry, the amount of substance (symbol ) in a given sample of matter is defined as a ratio () between the number of elementary entities () and the Avogadro constant (). The entities are usually molecules, atoms, ions, or ion pairs of a specified kind. The particular substance sampled may be specified using a subscript, e.g., the amount of sodium chloride (NaCl) would be denoted as . The unit of amount of substance in the International System of Units is the mole (symbol: mol), a base unit. Since 2019, the value of the Avogadro constant is defined to be exactly . Sometimes, the amount of substance is referred to as the chemical amount or, informally, as the "number of moles" in a given sample of matter.
Usage
Historically, the mole was defined as the amount of substance in 12 grams of the carbon-12 isotope. As a consequence, the mass of one mole of a chemical compound, in grams, is numerically equal (for all practical purposes) to the mass of one molecule or formula unit of the compound, in daltons, and the molar mass of an isotope in grams per mole is approximately equal to the mass number (historically exact for carbon-12 with a molar mass of 12 g/mol). For example, a molecule of water has a mass of about 18.015 daltons on average, whereas a mole of water (which contains water molecules) has a total mass of about 18.015 grams.
In chemistry, because of the law of multiple proportions, it is often much more convenient to work with amounts of substances (that is, number of moles or of molecules) than with masses (grams) or volumes (liters). For example, the chemical fact "1 molecule of oxygen () will react with 2 molecules of hydrogen () to make 2 molecules of water ()" can also be stated as "1 mole of will react with 2 moles of to form 2 moles of water". The same chemical fact, expressed in terms of masses, would be "32 g (1 mole) of oxygen will react with approximately 4.0304 g (2 moles of ) hydrogen to make approximately 36.0304 g (2 moles) of water" (and the numbers would depend on the isotopic composition of the reagents). In terms of volume, the numbers would depend on the pressure and temperature of the reagents and products. For the same reasons, the concentrations of reagents and products in solution are often specified in moles per liter, rather than grams per liter.
The amount of substance is also a convenient concept in thermodynamics. For example, the pressure of a certain quantity of a noble gas in a recipient of a given volume, at a given temperature, is directly related to the number of molecules in the gas (through the ideal gas law), not to its mass.
This technical sense of the term "amount of substance" should not be confused with the general sense of "amount" in the English language. The latter may refer to other measurements such as mass or volume, rather than the number of particles. There are proposals to replace "amount of substance" with more easily-distinguishable terms, such as enplethy and stoichiometric amount.
The IUPAC recommends that "amount of substance" should be used instead of "number of moles", just as the quantity mass should not be called "number of kilograms".
Nature of the particles
To avoid ambiguity, the nature of the particles should be specified in any measurement of the amount of substance: thus, a sample of 1 mol of molecules of oxygen () has a mass of about 32 grams, whereas a sample of 1 mol of atoms of oxygen () has a mass of about 16 grams.
Derived quantities
Molar quantities (per mole)
The quotient of some extensive physical quantity of a homogeneous sample by its amount of substance is an intensive property of the substance, usually named by the prefix "molar" or the suffix "per mole".
For example, the quotient of the mass of a sample by its amount of substance is its molar mass, for which the SI unit kilogram per mole or gram per mole may be used. This is about 18.015 g/mol for water, and 55.845 g/mol for iron. Similarly for volume, one gets the molar volume, which is about 18.069 millilitres per mole for liquid water and 7.092 mL/mol for iron at room temperature. From the heat capacity, one gets the molar heat capacity, which is about 75.385 J/(K⋅mol) for water and about 25.10 J/(K⋅mol) for iron.
Molar mass
The molar mass () of a substance is the ratio of the mass () of a sample of that substance to its amount of substance (): . The amount of substance is given as the number of moles in the sample. For most practical purposes, the numerical value of the molar mass in grams per mole is the same as that of the mean mass of one molecule or formula unit of the substance in daltons, as the mole was historically defined such that the molar mass constant was exactly 1 g/mol. Thus, given the molecular mass or formula mass in daltons, the same number in grams gives an amount very close to one mole of the substance. For example, the average molecular mass of water is about 18.015 Da and the molar mass of water is about 18.015 g/mol. This allows for accurate determination of the amount in moles of a substance by measuring its mass and dividing by the molar mass of the compound: . For example, 100 g of water is about 5.551 mol of water. Other methods of determining the amount of substance include the use of the molar volume or the measurement of electric charge.
The molar mass of a substance depends not only on its molecular formula, but also on the distribution of isotopes of each chemical element present in it. For example, the molar mass of calcium-40 is , whereas the molar mass of calcium-42 is , and of calcium with the normal isotopic mix is .
Amount (molar) concentration (moles per liter)
Another important derived quantity is the molar concentration () (also called amount of substance concentration, amount concentration, or substance concentration, especially in clinical chemistry), defined as the amount in moles () of a specific substance (solute in a solution or component of a mixture), divided by the volume () of the solution or mixture: .
The standard SI unit of this quantity is mol/m3, although more practical units are commonly used, such as mole per liter (mol/L, equivalent to mol/dm3). For example, the amount concentration of sodium chloride in ocean water is typically about 0.599 mol/L.
The denominator is the volume of the solution, not of the solvent. Thus, for example, one liter of standard vodka contains about 0.40 L of ethanol (315 g, 6.85 mol) and 0.60 L of water. The amount concentration of ethanol is therefore (6.85 mol of ethanol)/(1 L of vodka) = 6.85 mol/L, not (6.85 mol of ethanol)/(0.60 L of water), which would be 11.4 mol/L.
In chemistry, it is customary to read the unit "mol/L" as molar, and denote it by the symbol "M" (both following the numeric value). Thus, for example, each liter of a "0.5 molar" or "0.5 M" solution of urea () in water contains 0.5 moles of that molecule. By extension, the amount concentration is also commonly called the molarity of the substance of interest in the solution. However, as of May 2007, these terms and symbols are not condoned by IUPAC.
This quantity should not be confused with the mass concentration, which is the mass of the substance of interest divided by the volume of the solution (about 35 g/L for sodium chloride in ocean water).
Amount (molar) fraction (moles per mole)
Confusingly, the amount (molar) concentration should also be distinguished from the molar fraction (also called mole fraction or amount fraction) of a substance in a mixture (such as a solution), which is the number of moles of the compound in one sample of the mixture, divided by the total number of moles of all components. For example, if 20 g of is dissolved in 100 g of water, the amounts of the two substances in the solution will be (20 g)/(58.443 g/mol) = 0.34221 mol and (100 g)/(18.015 g/mol) = 5.5509 mol, respectively; and the molar fraction of will be .
In a mixture of gases, the partial pressure of each component is proportional to its molar fraction.
History
The alchemists, and especially the early metallurgists, probably had some notion of amount of substance, but there are no surviving records of any generalization of the idea beyond a set of recipes. In 1758, Mikhail Lomonosov questioned the idea that mass was the only measure of the quantity of matter, but he did so only in relation to his theories on gravitation. The development of the concept of amount of substance was coincidental with, and vital to, the birth of modern chemistry.
1777: Wenzel publishes Lessons on Affinity, in which he demonstrates that the proportions of the "base component" and the "acid component" (cation and anion in modern terminology) remain the same during reactions between two neutral salts.
1789: Lavoisier publishes Treatise of Elementary Chemistry, introducing the concept of a chemical element and clarifying the Law of conservation of mass for chemical reactions.
1792: Richter publishes the first volume of Stoichiometry or the Art of Measuring the Chemical Elements (publication of subsequent volumes continues until 1802). The term "stoichiometry" is used for the first time. The first tables of equivalent weights are published for acid–base reactions. Richter also notes that, for a given acid, the equivalent mass of the acid is proportional to the mass of oxygen in the base.
1794: Proust's Law of definite proportions generalizes the concept of equivalent weights to all types of chemical reaction, not simply acid–base reactions.
1805: Dalton publishes his first paper on modern atomic theory, including a "Table of the relative weights of the ultimate particles of gaseous and other bodies".
The concept of atoms raised the question of their weight. While many were skeptical about the reality of atoms, chemists quickly found atomic weights to be an invaluable tool in expressing stoichiometric relationships.
1808: Publication of Dalton's A New System of Chemical Philosophy, containing the first table of atomic weights (based on H = 1).
1809: Gay-Lussac's Law of combining volumes, stating an integer relationship between the volumes of reactants and products in the chemical reactions of gases.
1811: Avogadro hypothesizes that equal volumes of different gases (at same temperature and pressure) contain equal numbers of particles, now known as Avogadro's law.
1813/1814: Berzelius publishes the first of several tables of atomic weights based on the scale of m(O) = 100.
1815: Prout publishes his hypothesis that all atomic weights are integer multiple of the atomic weight of hydrogen. The hypothesis is later abandoned given the observed atomic weight of chlorine (approx. 35.5 relative to hydrogen).
1819: Dulong–Petit law relating the atomic weight of a solid element to its specific heat capacity.
1819: Mitscherlich's work on crystal isomorphism allows many chemical formulae to be clarified, resolving several ambiguities in the calculation of atomic weights.
1834: Clapeyron states the ideal gas law.
The ideal gas law was the first to be discovered of many relationships between the number of atoms or molecules in a system and other physical properties of the system, apart from its mass. However, this was not sufficient to convince all scientists of the existence of atoms and molecules, many considered it simply being a useful tool for calculation.
1834: Faraday states his Laws of electrolysis, in particular that "the chemical decomposing action of a current is constant for a constant quantity of electricity".
1856: Krönig derives the ideal gas law from kinetic theory. Clausius publishes an independent derivation the following year.
1860: The Karlsruhe Congress debates the relation between "physical molecules", "chemical molecules" and atoms, without reaching consensus.
1865: Loschmidt makes the first estimate of the size of gas molecules and hence of number of molecules in a given volume of gas, now known as the Loschmidt constant.
1886: van't Hoff demonstrates the similarities in behaviour between dilute solutions and ideal gases.
1886: Eugen Goldstein observes discrete particle rays in gas discharges, laying the foundation of mass spectrometry, a tool subsequently used to establish the masses of atoms and molecules.
1887: Arrhenius describes the dissociation of electrolyte in solution, resolving one of the problems in the study of colligative properties.
1893: First recorded use of the term mole to describe a unit of amount of substance by Ostwald in a university textbook.
1897: First recorded use of the term mole in English.
By the turn of the twentieth century, the concept of atomic and molecular entities was generally accepted, but many questions remained, not least the size of atoms and their number in a given sample. The concurrent development of mass spectrometry, starting in 1886, supported the concept of atomic and molecular mass and provided a tool of direct relative measurement.
1905: Einstein's paper on Brownian motion dispels any last doubts on the physical reality of atoms, and opens the way for an accurate determination of their mass.
1909: Perrin coins the name Avogadro constant and estimates its value.
1913: Discovery of isotopes of non-radioactive elements by Soddy and Thomson.
1914: Richards receives the Nobel Prize in Chemistry for "his determinations of the atomic weight of a large number of elements".
1920: Aston proposes the whole number rule, an updated version of Prout's hypothesis.
1921: Soddy receives the Nobel Prize in Chemistry "for his work on the chemistry of radioactive substances and investigations into isotopes".
1922: Aston receives the Nobel Prize in Chemistry "for his discovery of isotopes in a large number of non-radioactive elements, and for his whole-number rule".
1926: Perrin receives the Nobel Prize in Physics, in part for his work in measuring the Avogadro constant.
1959/1960: Unified atomic mass unit scale based on m(C) = 12 u adopted by IUPAP and IUPAC.
1968: The mole is recommended for inclusion in the International System of Units (SI) by the International Committee for Weights and Measures (CIPM).
1972: The mole is approved as the SI base unit of amount of substance.
2019: The mole is redefined in the SI as "the amount of substance of a system that contains specified elementary entities".
| Physical sciences | Chemistry: General | null |
8082793 | https://en.wikipedia.org/wiki/Breakdown%20%28vehicle%29 | Breakdown (vehicle) | A vehicle breakdown is a mechanical or electrical failure of a motor vehicle in such a way that the underlying problem prevents the vehicle from being operated or impedes the vehicle's operation so significantly that it is very difficult, nearly impossible, or else dangerous to operate.
Vehicle breakdowns have various causes. Depending on the severity, the vehicle may need to be towed to an automobile repair shop or fixed on-site by roadside assistance or a mobile mechanic. With other problems, the driver may be able to operate the vehicle seemingly normally for some time, but the vehicle will need an eventual repair. Many vehicle owners with personal economic difficulty or a busy schedule may wait longer than they should to get necessary maintenance or repairs made to their vehicles, thereby increasing their chances of a breakdown, inducing further damage to the vehicle, or else causing more danger.
Severity
There are various levels of a vehicle's disability:
Total Breakdown
A total breakdown is when the vehicle becomes totally immobile and cannot be driven even a short distance, thereby necessitating a tow to an auto mechanic, or an on-site repair from roadside assistance or a mobile technician.
Partial Breakdown
In a partial breakdown, the vehicle may still be operable, but its operation may become more limited or more dangerous, or else its continued operation may contribute to further damage to the vehicle. When this occurs, it may be possible to drive the vehicle to a repair shop or a safer location for a mobile repair, thereby avoiding a tow.
Causes
In 2014, The Royal Automobile Club (RAC) attended almost two million breakdowns in the United Kingdom. Battery problems were the most common cause of a car breakdown, accounting for more than 450,000 call-outs.
Top 10 Causes
Dead or faulty battery
Engine issues
Faulty alternator
Damaged tyre & wheel
Electrical issues
Starter motor problem
Damaged clutch wire
Brake problems
Fuel problems
Lost car keys
Source:
Contributing Factors
Roadside assistance data collected, analyzed and published by AAA provides the following statistical insights into vehicle breakdowns in the United States and Canada:
Age
Vehicles ten years and older are twice as likely to end up immobilized on the side of the road compared to newer vehicles, and the odds of needing a tow quadruples.
Vehicles fewer than five years old have a higher proportion of tire, key and fuel-related issues than older vehicles.
One in five service calls for a newer vehicle required a tow to a repair facility.
Vehicles between six and ten years old have the highest proportion of battery-related issues, as most batteries have a three-to-five-year life.
Maintenance
Millions of roadside breakdowns each year could be prevented with basic vehicle maintenance.
35 percent of Americans have skipped or delayed service or repairs that were recommended by a mechanic or specified by the factory maintenance schedule.
Region
Drivers in the West experienced the most breakdowns, followed by the South, the Northeast and the Midwest.
Season
Roadside assistance calls peak in the summer (8.3 million) followed by winter (8.1 million), fall (7.8 million) and spring (7.7 million).
Coverage
When a breakdown occurs, the motorist may be able to have the tow and/or repair covered by a third party:
Factory Warranty
A vehicle manufacturer's warranty often covers the repair and subsequent towing expense if the root cause of the repair is deemed warrantable and is repaired at a certified dealership.
Extended Warranty/Service Contract
A service contract is a type of extended warranty (although may not be legally considered a true warranty) that is typically available for purchase at the time of the vehicle sale and usually includes roadside assistance.
Insurance
An auto insurance policy provides liability coverage for any damage caused by a vehicle, however, some comprehensive policies also cover repairs, roadside assistance services, tow and rental expenses.
Roadside Assistance
A roadside assistance plan may be available for purchase through club membership or in some cases, as an add-on service from another subscription, such as a mobile phone.
| Technology | Concepts of ground transport | null |
1219865 | https://en.wikipedia.org/wiki/Panthera%20spelaea | Panthera spelaea | Panthera spelaea, commonly known as the cave lion (or less commonly as the steppe lion), is an extinct Panthera species that was native to Eurasia and northwest North America during the Pleistocene epoch. Genetic analysis of ancient DNA has revealed that while closely related, it was a distinct species genetically isolated from the modern lion (Panthera leo), with the genetic divergence between the two species estimated at around 500,000 years ago. The earliest fossils of the P. spelaea lineage (either regarded as the separate species Panthera fossilis or the subspecies P. spelaea fossilis) in Eurasia date to around 700,000 years ago (with possible late Early Pleistocene records). It is closely related and probably ancestral to the American lion (Panthera atrox). The species ranged from Western Europe to eastern Beringia in North America, and was a prominent member of the mammoth steppe fauna, and an important apex predator across its range. It became extinct about 13,000 years ago. It closely resembled living lions with a coat of yellowish-grey fur though unlike extant lions, males appear to have lacked manes.
Panthera spelaea interacted with both Neanderthals and modern humans, who used their pelts and in the case of the latter, depicted them in artistic works.
Taxonomy
In 1774, Zoolithenhöhle cave near the village of Burggaillenreuth in Bavaria, southern Germany was brought to scientific attention by Johan Friedrich Esper, who realised that the bones of extinct animals were present in the cave. In 1810, a fossil skull from the cave was given the scientific name Felis spelaea by Georg August Goldfuss. It possibly dates to the Last Glacial Period.
Several anatomical studies of remains of Panthera spelaea were conducted during the early-mid 19th century, who generally agreed that the species had lion affinities. During the 20th century and the first decade of the 21st century, Panthera spelaea was often regarded as a subspecies of the modern lion, and therefore as Panthera leo spelaea. One author considered the cave lion to be more closely related to the tiger based on a comparison of skull shapes, and proposed the scientific name Panthera tigris spelaea.
Results from morphological studies showed that it is distinct in cranial and dental anatomy to justify the specific status of Panthera spelaea. Results of phylogenetic studies also support this assessment.
In 2001, the subspecies Panthera spelaea vereshchagini was proposed for seven specimens found in Siberia and Yukon, which have smaller skulls and teeth than the average P. spelaea. Before 2020, genetic analysis using ancient DNA provided no evidence for their distinct subspecific status; DNA signatures from P. spelaea from Europe and Alaska were indistinguishable, suggesting one large panmictic population. However, analysis of mitochondrial genome sequences from 31 cave lions showed that they fall into two monophyletic clades. One lived across western Europe and the other was restricted to Beringia during the Pleistocene. For this reason, the Beringian population is considered a distinct subspecies, P. s. vereshchagini.
Evolution
Lion-like pantherine felids first appeared in the Tanzanian Olduvai Gorge about . These cats dispersed into Eurasia from East Africa around the end of the Early Pleistocene and the beginning of the Middle Pleistocene, giving rise to Panthera fossilis. The oldest widely accepted fossils of P. fossilis in Europe date to around 700,000 years ago, such as that from Pakefield in England, with possible older fossils from Western Siberia dating to the late Early Pleistocene, with a 2024 study suggesting a presence in Spain by 1 million years ago. Different authors considered Panthera fossils as either a distinct species ancestral to P. spelaea, or as a subspecies of P. spelaea. Recent nuclear genomic evidence suggest that interbreeding between modern lions and all Eurasian fossil lions took place up until 500,000 years ago, but by 470,000 years ago, no subsequent interbreeding between the two lineages occurred.
The following cladogram shows the genetic relationship between P. spelaea and other pantherine cats.The arrival of Panthera (spelaea) fossilis in Europe was part of a faunal turnover event around the Early-Middle Pleistocene transition in which many of the species that characterised the preceeding late Villafranchian became extinct. In the carnivore guild, this notably included the giant hyena Pachycrocuta and the sabertooth cat Megantereon. Following the arrival of Panthera (spelaea) fossilis the lion-sized sabertooth cat Homotherium and the "European jaguar" Panthera gombaszoegensis became much rarer, ultimately becoming extinct in Europe during the late Middle Pleistocene, with competition with lions suggested to be a likely important factor.
Specimens intermediate between P. fossilis and Late Pleistocene P. spelaea are referred to as the subspecies P. s. intermedia. The transition from P. fossilis to Late Pleistocene P. spelaea shows significant reduction in body size, as well as changes in skull and tooth morphology. Mitochondrial DNA sequence data from fossil lion remains show that the American lion represents a sister group of Late Pleistocene P. spelaea, and likely arose when an early P. spelaea population became isolated south of the Cordilleran Ice Sheet. Initially this was suggested to be around 340,000 years ago, but later studies suggested that the split between the two species was probably younger, around 165,000 years ago, consistent with the late first appearance of P. spelaea in Eastern Beringia (now Alaska and adjacent regions) during the Illinoian (around 190-130,000 years ago).
Characteristics
Carvings and cave paintings of cave lions, which were discovered in the Lascaux and Chauvet Caves in France, were dated to 15,000 to 17,000 years old. A drawing in the Chauvet cave depicts two cave lions walking together. The one in the foreground is slightly smaller than the one in the background, which has been drawn with a scrotum and without a mane. Such cave paintings suggest that male cave lions completely lacked manes, or at most had very small manes.
Early members of the cave lion lineage assigned to Panthera (spelaea) fossilis during the Middle Pleistocene were considerably larger than individuals of P. spelaea from the Last Glacial Period and modern lions, with some of these individuals having an estimated length of , shoulder height of and body mass of , respectively, making them among the largest cats to have ever lived. The Late Pleistocene Panthera spelaea spelaea was noticeably smaller though still large relative to living cats, with an estimated length of and shoulder height of , respectively, The species showed a progressive size reduction over the course of the Last Glacial Period up until its extinction, with the last P. spelaea populations comparable in size to small-sized modern lions, with a body mass of only , a body length of and shoulder height of respectively.
P. spelaea had a relatively longer and narrower muzzle compared to that of the extant lion, with the zygomatic region being strongly arched, with the carnassial teeth having differences in cusp morphology (displaying preparastyles). Like modern lions, females were smaller than males. Compared to the earlier P. (spelaea) fossilis, Late Pleistocene P. spelaea spelaea differs (in addition to previously mentioned size differences) in having larger incisor teeth, more narrow and flattened canines, as well narrower upper and lower third and fourth premolars, which display some differences in cusp morphology, with the lower first molar being narrower and more elongate. The orbits (eye sockets) of P. spelaea spelaea are also relatively larger and muzzle marginally narrower compared to P. (spelaea) fossilis, with the nasal region also being proportionally narrower, while the postorbital and mastoid regions of the skull are wider, with the tympanic bullae being more inflated.
In 2016, hair found near the Maly Anyuy River was identified as cave lion hair through DNA analysis. Comparison with hair of a modern lion revealed that cave lion hair was probably similar in colour as that of the modern lion, though slightly lighter. In addition, the cave lion is thought to have had a very thick and dense undercoat comprising closed and compressed yellowish-to-white wavy downy hair with a smaller mass of darker-coloured guard hairs, possibly an adaptation to the Ice Age climate. While juveniles fur coat colour was yellowish, adult cave lions are suggested to have had grey fur.
Distribution and habitat
During the Last Glacial Period, P. spelaea formed a contiguous population across the mammoth steppe, from Western Europe to northwest North America. It was widely distributed in the Iberian Peninsula, Italian Peninsula, Southeast Europe, Great Britain, Central Europe, the East European Plain, the Ural Mountains, most of Northeast Asia (ranging as far south as Northeast China and possibly the Korean peninsula), and across the Bering land bridge into Alaska and Yukon. The cave lion had a wide elevation range, with finds extending up over above sea level in the European Alps and in Buryatia in Northern Asia, though they probably did not occupy mountainous habitats all-year round. The cave lion probably inhabited predominantly open habitats such as steppe and grasslands although it would have also have occurred in open woodlands as well. While during the Last Glacial Period it was often associated with cold environments, the species also inhabited temperate environments, such as in Europe during the Last Interglacial/Eemian.
Paleobiology
Ecology
P. spelaea was one of the keystone species of the mammoth steppe, being one of the main apex predators alongside the gray wolf, cave hyena and brown bear. Large amounts of bones belonging to P. spelaea were excavated in caves, where bones of cave hyena, cave bear and Paleolithic artefacts were also found. Despite their common name, "cave lions" probably only infrequently if ever used caves, and were present in regions where caves were absent. Some of these accumulations of cave lion bones in cave hyena dens have been attributed to confrontations between cave hyenas and cave lions over carcasses, with the remains of cave lions killed in these confrontations subsequently transported to the dens.
Isotopic analyses of bone collagen samples extracted from remains in Europe and East Beringia indicate that reindeer were particularly prominent in the diet of cave lions in these regions during the Last Glacial Period. Cave lions also seem to have opportunistically preyed on the cubs of cave bears. Isotopic analysis of other European specimens suggests a diet including wild horse, woolly mammoth and cave bears for these individuals. It may have sought out hibernating bears in montane caves as a food source during the winter. Bite marks found on the bones of straight-tusked elephants in Neumark Nord, Germany, dating to the Last Interglacial, have been suggested to be the result of scavenging by cave lions.
Other possible prey species were giant deer, red deer, muskox, aurochs, wisent, steppe bison, and young woolly rhinoceros. It likely competed for prey with the European leopards, cave hyenas, brown bears and grey wolves in Eurasia, along with short-faced bears, Homotherium, and Beringian wolves in Beringia.
Social behavior
Whether or not cave lions existed in prides like modern lions is unclear. Isotopic analysis on cave lions by Hervé Bocherens and colleagues lead them to suggest that cave lions may have been solitary, due to cave lions shifting their diets after the disappearance of cave hyenas, carcasses being consumed the cave hyenas as well, suggests they were at a competitive disadvantage, and the scattering of isotopic data between individuals. Some other authors have also argued that the absence of manes in cave lions suggests that cave lions did not live in prides, given the importance of manes in the social hierarchy of modern lions. Boeskorov et al. 2021 suggested both European and Beringian cave lions may have hunted in larger prides than modern lions because sexual dimorphism in cave lions was more pronounced than in modern African lions and solitary big cats. However, they admitted the data is insufficient to come down to a certain conclusion.
Cave lion cubs appear to have lived in dens during their earliest stages of life, like modern lion cubs and were likely solely raised by females, like living Panthera species.
Relationship with humans
Cave lions were hunted and their pelts exploited in Europe by Neanderthals during the Middle Paleolithic, and during the Upper Paleolithic by modern humans in Spain as evidenced in the La Garma site dating to the Magdalenian. Modern humans also drew cave paintings of cave lions, engraved their likeness on bones and created sculptures of them, including the famous anthropomorphic lion-man figure from Hohlenstein-Stadel cave in Germany dating to around 41-35,000 years ago with the body of a human and the head of a lion. Cave lion canines with perforated holes may have been worn as personal ornaments. Decorated stones with engravings representing cave lions have been found in southern Italy.
Extinction
Radiocarbon dating suggests that the species went extinct approximately simultaneously across its range during the last few thousand years of the Late Pleistocene, around 14-15,000 years ago, possibly surviving around 1000 years later in the far east North American portion of its range. This timing roughly corresponds to the onset of the Bølling–Allerød Interstadial warm period and the consequent collapse of the mammoth steppe ecosystem. The precise cause of its extinction is unclear, but possibly involved environmental change from open habitats to closed forests, changes in prey abundance, as well as human impact, though it is difficult to distentangle the precise causes of its extinction. Cave lions appear to have undergone a population bottleneck that considerably reduced their genetic diversity between 47,000 and 18,000 years ago, probably driven at least in part by climatic instability.
Mummified specimens
In 2008, a well-preserved mature cave lion specimen was unearthed near the Maly Anyuy River in Chukotka Autonomous Okrug, which still retained some clumps of hair.
In 2015, two frozen cave lion cubs, estimated to be between 25,000 and 55,000 years old, were discovered close to the Uyandina River in Yakutia, Siberia in permafrost.
Research results indicate that the cubs were likely barely a week old at the time of their deaths, as their milk teeth had not fully erupted. Further evidence suggests the cubs were hidden at a den site until they were strong enough to follow their mother back to the pride, as with modern lions. Researchers believe that the cubs were trapped and killed by a landslide, and that the absence of oxygen underground hindered their decomposition and allowed the cubs to be preserved in such good condition. A second expedition to the site where the cubs were found was planned for 2016, in hopes of finding either the remains of a third cub or possibly the cubs' mother.
In 2017, another frozen specimen, thought to be a lion cub, was found in Yakutia on the banks of the Tirekhtyakh River (), a tributary of the Indigirka River. This male cub was thought to be slightly older than the 2015 cubs at the time of its death; it is estimated to have been around one and a half to two months. In 2018, another preserved carcass of a cub was found in a location away. It was considered to be around a month old when it died approximately 50,000 years ago, and presumed to be a sibling of the male cub. However, carbon dating showed them to have lived about 15,000 years apart, with the female estimated to have lived 28,000 years ago, and the male 43,448 years ago. Both cubs were well preserved, albeit with a few damages, with the female possibly being the "best preserved" animal discovered from the Ice age.
| Biology and health sciences | Felines | Animals |
1220214 | https://en.wikipedia.org/wiki/American%20lion | American lion | The American lion (Panthera atrox (), with the species name meaning "savage" or "cruel", also called the North American lion) is an extinct pantherine cat native to North America during the Late Pleistocene from around 130,000 to 12,800 years ago. Genetic evidence suggests that its closest living relative is the lion (Panthera leo), with the American lion representing an offshoot from the lineage of the largely Eurasian cave lion (Panthera spelaea), from which it is suggested to have split around 165,000 years ago. Its fossils have been found across North America, from Canada to Mexico. It was about 25% larger than the modern lion, making it one of the largest known felids to ever exist, and an important apex predator.
The American lion became extinct as part of the end-Pleistocene extinction event along with most other large animals across the Americas. The extinctions followed human arrival in the Americas. Proposed factors in its extinction include climatic change reducing viable habitat, as well as human hunting of herbivore prey causing a trophic cascade.
History and taxonomy
Initial discovery and North American fossils
The first specimen now assigned to Panthera atrox was collected in the 1830s by William Henry Huntington, Esq., who announced his discovery to the American Philosophical Society on April 1, 1836 and placed it with other fossils from Huntington's collection in the Academy of Natural Sciences in Philadelphia. The specimen had been collected in ravines in Natchez, Mississippi that were dated to the Pleistocene; the specimen consisted only of a partial left mandible with 3 molars and a partial canine. The fossils did not get a proper description until 1853 when Joseph Leidy named the fragmentary specimen (ANSP 12546) Felis atrox ("savage cat"). Leidy named another species in 1873, Felis imperialis, based on a mandible fragment from Pleistocene gravels in Livermore Valley, California. F. imperialis however is considered a junior synonym of Panthera atrox. A replica of the jaw of the first American lion specimen to be discovered can be seen in the hand of a statue of famous paleontologist Joseph Leidy, currently standing outside the Academy of Natural Sciences in Philadelphia.
Few additional discoveries came until 1907, when the American Museum of Natural History and College, Alaska collected several Panthera atrox skulls in a locality originally found in 1803 by gold miners in Kotzebue, Alaska. The skulls were referred to a new subspecies of Felis (Panthera) atrox in 1930, Felis atrox "alaskensis". Despite this, the species did not get a proper description and is now seen as a nomen nudum synonymous with Panthera atrox. Further south in Rancho La Brea, California, a large felid skull was excavated and later described in 1909 by John C. Merriam, who referred it to a new subspecies of Felis atrox, Felis atrox bebbi. The subspecies is synonymous with Panthera atrox.
Throughout the early to mid 1900s, dozens of fossils of Panthera atrox were excavated at La Brea, including many postcranial elements and associated skeletons. The fossils were described by Merriam & Stock in detail in 1932, who synonymized many previously named taxa with Felis atrox. At least 80 individuals are known from La Brea Tar Pits and the fossils define the subspecies, giving a comprehensive view of the taxon. It was not until 1941 that George Simpson moved Felis atrox to Panthera, believing that it was a subspecies of jaguar. Simpson also referred several fossils from central Mexico, even as far south as Chiapas, as well as Nebraska and other regions of the western US, to P. atrox. 1971 witnessed the description of fragmentary remains from Alberta, Canada that extended P. atroxs range north. In 2009, an entrapment site at Natural Trap Cave, Wyoming was briefly described and is the second most productive Panthera atrox-bearing fossil site. It most importantly contains well-preserved mitochondrial DNA of many partial skeletons.
Panthera onca mesembrina and possible South American material
In the 1890s in the "Cueva del Milodon" in southern Chile, fossil collector Rodolfo Hauthal collected a fragmentary postcranial skeleton of a large felid that he sent to Santiago Roth. Roth described them as a new genus and species of felid, "Iemish listai" in 1899. However, the name is considered a nomen nudum. In 1904, Roth reassessed the phylogenetic affinities of "Iemish" and named it Felis listai and referred several cranial and fragmentary postcranial elements to the taxon. Notably, several mandibles, a partial skull, and pieces of skin were some of the specimens referred. In 1934, Felis onca mesembrina was named by Angel Cabrera based on that partial skull from "Cueva del Milodon" and the other material from the site was referred to it. The skull (MLP 10-90) was lost, and was only illustrated by Cabrera.
Further material, including feces and mandibles, was referred to as F. onca mesembrina from Tierra del Fuego, Argentina and other southern sites in Chile.
In 2016, the subspecies was referred to Panthera onca in a genetic study, which supported its identity as a subspecies of jaguar. Later in 2017, one study synoymised P. onca mesembrina with Panthera atrox based on morphological similarities, though this does not have broad acceptance.
Evolution
The American lion was initially considered a distinct species of Pantherinae, and designated as Panthera atrox, which means "cruel" or "fearsome panther" in Latin. Some paleontologists accepted this view, but others considered it to be a type of lion closely related to the modern lion (Panthera leo) and its extinct relative, the Eurasian cave lion (Panthera leo spelaea or P. spelaea). It was later assigned as a subspecies of P. leo (P. leo atrox) rather than as a separate species. Most recently, both spelaea and atrox have been treated as full species.
Cladistic studies using morphological characteristics have been unable to resolve the phylogenetic position of the American lion. One study considered the American lion, along with the cave lion, to be most closely related to the tiger (Panthera tigris), citing a comparison of the skull; the braincase, in particular, appears to be especially similar to the braincase of a tiger. Another study suggested that the American lion and the Eurasian cave lion were successive offshoots of a lineage leading to a clade which includes modern leopards and lions. A more recent study comparing the skull and jaw of the American lion with other pantherines concluded that it was not a lion but a distinct species. It was proposed that it arose from pantherines that migrated to North America during the mid-Pleistocene and gave rise to American lions and jaguars (Panthera onca). Another study grouped the American lion with P. leo and P. tigris, and ascribed morphological similarities to P. onca to convergent evolution, rather than phylogenetic affinity.
Mitochondrial DNA sequence data from fossil remains suggests that the American lion (P. atrox) represents a sister lineage to Late Pleistocene populations of the Eurasian cave lion (P. spelaea), and likely arose when an early cave lion population became isolated south of the North American continental ice sheet. While initial studies suggested that the divergence between American and Eurasian cave lions took place around 340,000 years ago, later studies suggested that the split took place considerably later, around 165,000 years ago, consistent with the earliest appearance of cave lions in eastern Beringia (now Alaska) during the Illinoian (190-130,000 years ago).
Genetic studies indicate that the living lion is the closest living relative of P. atrox and P. spelaea. Genome-wide sequencing of modern lions and Eurasian cave lions suggests that the lineage of the cave lion and American lion diverged from that of the modern lion around 500,000 years ago.
Description
The American lion is estimated to have measured from the tip of the nose to the base of the tail and stood at the shoulder. Panthera atrox was sexually dimorphic, with an approximate range of between 235kg to 523 kg (518lbs-1153lbs) in males and 175kg to 365 kg (385lbs-805lbs) for females. A separate study found American lions were more sexually dimorphic than modern lions in terms of size: American lion males being 1.4 times larger than females, compared to modern male lions being 1.26 times larger. In 2008, the American lion was estimated to weigh up to . A study in 2009 showed an average weight of for males and for the largest specimen analyzed.
Panthera atrox had limb bones more robust than those of an African lion, and comparable in robustness to the bones of a brown bear; also its limbs were 10 % longer than extant African lion in relation to skull length. About 80 American lion individuals have been recovered from the La Brea Tar Pits in Los Angeles, so their morphology is well known. Their features strongly resemble those of modern lions, but they were considerably larger, similar to P. spelaea and the Pleistocene Natodomeri lion of eastern Africa.
Preserved skin remains found with skeletal material considered by some to belong to the American lion found in caves in Patagonia is reddish in colour, though the attribution of Patagonian Panthera remains to P. atrox is highly controversial and not accepted by many authors. Preserved fur of the closely related P. spelaea found in Siberia is yellowish in colour, with cave art of European P. spelaea indicating that males lacked substantial manes unlike modern lions. These characteristics may also apply to P. atrox.
Distribution
The earliest lions known in the Americas south of Alaska are from the Sangamonian Stage (equivalent to the global Last Interglacial ~130-115,000 years ago) during which American lions rapidly dispersed across North America, with their distribution ultimately ranging from Canada to southern Mexico and from California to the Atlantic coast. It was generally not found in the same areas as the jaguar, which favored forests over open habitats. It was absent from eastern Canada and the northeastern United States, perhaps due to the presence of dense boreal forests in the region. Farther south, fossilised remains of the American lion have been discovered in Extinction Cave, Belize. The American lion was formerly believed to have colonized northwestern South America as part of the Great American Interchange. However, the fossil remains found in the tar pits of Talara, Peru actually belong to an unusually large jaguar. On the other hand, fossils of a large felid from late Pleistocene localities in southern Chile and Argentina traditionally identified as an extinct subspecies of jaguar, Panthera onca mesembrina, have been considered by some authors actually represent remains of the American lion, though this interpretation is highly controversial, with many authors favouring a jaguar attribution for these remains.
Paleobiology
The American lion inhabited savannas and grasslands like the modern lion.
Predatory behavior
American lions likely preyed on deer, horses, camels, tapirs, American bison, mammoths, and other large ungulates (hoofed mammals). Paired nitrogen and carbon isotopic evidence from Natural Trap Cave in Wyoming reveals that the extant pronghorn was an important food source for American lions, which probably hunted them regularly, although probably also could be due to kleptoparasitism from the kills of Miracinonyx (sometimes called the "American cheetah").
Analyses of dental microwear suggest that the American lion actively avoided bone just like the modern cheetah (more so than Smilodon). Panthera atrox has the highest proportion of canine breakage in La Brea, suggesting a consistent preference for larger prey than contemporary carnivores. Dental microwear additionally suggests that carcass utilization slightly declined over time (~30,000 BP to 11,000 radiocarbon BP) in Panthera atrox. The fragment of a femur from a gray wolf from the La Brea Tar Pits shows evidence of a violent bite which possibly amputated the leg. Researchers believe that Panthera atrox is a prime candidate for the injury, due to its bite force and bone shearing ability. Based on skull width, it is estimated that a 347 kilogram American lion would have a bite force of 2,830 newtons.
Social behavior
Whether American lion formed prides like modern lions or lived solitary lives like tigers is unknown. American lions likely descended from Panthera spelaea, which was likely a solitary animal, based on fossil evidence and several isotopic studies. However, if this would apply to the American Lion is unclear.
One study suggests American lion probably lived in prides like modern lions due to the large amounts of young males at dispersal age and low number of young females found at the tar pits. They argued that female American lions were less likely to end up in the tar pits because they were more likely to remain in their natural prides. However, the authors of the paper admit the higher ratio of males might just because the males in felids are more likely to disperse and thus fall into such traps, and that there's ultimately not any true evidence for or against sociality in American lions due a small sample size.
The remains of American lions are not as abundant as those of other predators like Smilodon fatalis or dire wolves (Aenocyon dirus) at the La Brea Tar Pits. This suggests that they were better at evading entrapment, possibly due to greater intelligence. Despite its rarity, the high ratio of juveniles to adults recovered at the tar pits would suggest possible gregariousness in Panthera atrox. But it’s rarity in the tar pits would suggests that it was possibly more solitary than Smilodon and Aenocyon or was gregarious but lived in low densities similar to African wild dogs.
Extinction
The American lion went extinct as part of the end-Pleistocene extinctions around 13-12,000 years ago, approximately simultaneously with most large (megafaunal) mammals across the Americas. The most recent fossil, from Edmonton, Canada dates to ~12,877 calibrated years Before Present, and is 400 years younger than the youngest cave lion in Alaska. These extinctions post-date human arrival to the Americas. The causes of the extinctions have been long the subject of controversy, with most authors positing climate change, humans or some combination of the two as the causes of the extinctions. A 2017 study suggested that the viable habitat for the American lion in North America had been greatly reduced over the course of the Last Glacial Period, which would have made it more vulnerable to extinction. Other authors have suggested that the extinction of the American lion and other competing carnivores like dire wolves, and the sabertooth cats Smilodon and Homotherium may have been due to trophic cascade effects caused by Paleoindian hunting of herbivores. These authors suggested that the herbivores already probably existed at low population numbers prior to Paleoindian arrival due to their abundance being limited by predators, rather than being at the carrying capacity of the ecosystem based on food resources. Due to humans having a more flexible omnivorous diet they may have been less subject to competition with other apex predators, allowing their population numbers to increase even as the number of herbivores declined.
| Biology and health sciences | Felines | Animals |
1221281 | https://en.wikipedia.org/wiki/Diapause | Diapause | In animal dormancy, diapause is the delay in development in response to regular and recurring periods of adverse environmental conditions. It is a physiological state with very specific initiating and inhibiting conditions. The mechanism is a means of surviving predictable, unfavorable environmental conditions, such as temperature extremes, drought, or reduced food availability. Diapause is observed in all the life stages of arthropods, especially insects.
Activity levels of diapausing stages can vary considerably among species. Diapause may occur in a completely immobile stage, such as the pupae and eggs, or it may occur in very active stages that undergo extensive migrations, such as the adult monarch butterfly, Danaus plexippus. In cases where the insect remains active, feeding is reduced and reproductive development is slowed or halted.
Embryonic diapause, a somewhat similar phenomenon, occurs in over 130 species of mammals, possibly even in humans, and in the embryos of many of the oviparous species of fish in the order Cyprinodontiformes.
Phases of insect diapause
Diapause in insects is a dynamic process consisting of several distinct phases. While diapause varies considerably from one taxon of insects to another, these phases can be characterized by particular sets of metabolic processes and responsiveness of the insect to certain environmental stimuli. For example, Sepsis cynipsea flies primarily use temperature to determine when to enter diapause. Similarly, Chrysoperla plorabunda lacewings regulate their reproductive cycle using daylight length, with adults entering reproductive diapause when there are less than 12-13 hours of daylight. Diapause can occur during any stage of development in arthropods, but each species exhibits diapause in specific phases of development. Reduced oxygen consumption is typical as is reduced movement and feeding. In Polistes exclamans, a social wasp, only the queen is said to be able to undergo diapause.
Comparison of diapause periods
The sensitive stage is the period when stimulus must occur to trigger diapause in the organism. Examples of sensitive stage/diapause periods in various insects:
{| class="wikitable"
|-
| Scientific name || Common name || Sensitive stage || Diapause
|-
|Diatraea grandiosella || Southwestern corn borer || early larval || late larval
|-
|Sarcophaga crassipalpis || Flesh fly || early larval || pupa
|-
|Sarcophaga argyrostoma || Flesh fly || mid to late larval || pupa
|-
|Manduca sexta || Tobacco hornworm || late embryonic (egg) to late larval || pupa
|-
|Leptinotarsa decemlineata || Colorado potato beetle || early adult || late adult
|-
|Bombyx mori || Silkworm || late embryonic (egg) to early larval || embryonic
|-
|Lymantria dispar || Spongy moth || late embryonic || late embryonic
|-
|Danaus plexippus || Monarch butterfly || early adulthood || adulthood
|-
|Acronicta rumicis
|Knott grass moth
|mid larval
|mid larval
|-
|Cydia pomonella
|Codling moth
|early to mid larval
|mid larval
|-
|Gynaephora groenlandica
|Arctic woolly bear moth
|mid larval
|mid larval
|-
|Cuterebra fontinella || Mouse botfly || mid larval || pupa
|-
|Nothobranchius furzeri
|turquoise killifish
|egg
|egg
|}
Induction
The induction phase occurs at a genetically predetermined stage of life, and occurs well in advance of the environmental stress. This sensitive stage may occur within the lifetime of the diapausing individual, or in preceding generations, particularly in egg diapause. During this phase, insects are responsive to external cues called token stimuli, which trigger the switch from direct development pathways to diapause pathways. Token stimuli can consist of changes in photoperiod, thermoperiod, or allelochemicals from food plants. These stimuli are not in themselves favourable or unfavourable to development, but they herald an impending change in environmental conditions.
Preparation
The preparation phase usually follows the induction phase, though insects may go directly from induction to initiation without a preparation phase. During this phase, insects accumulate and store molecules such as lipids, proteins, and carbohydrates. These molecules are used to maintain the insect throughout diapause and to provide fuel for development following diapause termination. Composition of the cuticle may be altered by changing hydrocarbon composition and by adding lipids to reduce water loss, making the organism resistant to desiccation. Diapausing puparia of the flesh fly, Sarcophaga crassipalpis, increase the amount of cuticular hydrocarbons lining the puparium, effectively reducing the ability of water to cross the cuticle.
Initiation
Photoperiod is the most important stimulus initiating diapause. The initiation phase begins when morphological development ceases. In some cases, this change may be very distinct and can involve moulting into a specific diapause stage, or be accompanied by color change. Enzymatic changes may take place in preparation for cold hardening. For example, only diapausing adults of the fire bug, Pyrrhocoris apterus, have the enzymatic complement that allows them to accumulate polyhydric alcohols, molecules that help to lower their freezing points and thus avoid freezing. Insects may also undergo behavioural changes and begin to aggregate, migrate, or search for suitable overwintering sites.
Maintenance
During the maintenance phase, insects experience lowered metabolism and developmental arrest is maintained. Sensitivity to certain stimuli which act to prevent termination of diapause, such as photoperiod and temperature, is increased. At this stage, insects are unresponsive to changes in the environment that will eventually trigger the end of diapause, but they grow more sensitive to these stimuli as time progresses.
Termination
In insects that undergo obligate diapause, termination may occur spontaneously, without any external stimuli. In facultative diapausers, token stimuli must occur to terminate diapause. These stimuli may include chilling, freezing, or contact with water, depending on the environmental conditions being avoided. These stimuli are important in preventing the insect from terminating diapause too soon, for instance in response to warm weather in late fall. In the Edith's checkerspot butterfly, individuals must receive enough sunlight in order to terminate the diapause stage and become a fully grown butterfly. Termination may occur at the height of unfavourable conditions, such as in the middle of winter. Over time, depth of diapause slowly decreases until direct development can resume, if conditions are favourable. Termination can also occur in specific time frames linked to reproductive periods, such as in the beetle Colaphellus bowringi: diapause ends for spring-reproducing beetles between late February and early April and for autumn-reproducing beetles between mid August and early October.Post-diapause quiescence
Diapause frequently ends prior to the end of unfavourable conditions and is followed by a state of quiescence from which the insect can arouse and begin direct development, should conditions change to become more favourable. This allows the insect to continue to withstand harsh conditions while being ready to take advantage of good conditions as soon as possible.
Regulation
Diapause in insects is regulated at several levels. Environmental stimuli interact with genetic pre-programming to affect neuronal signalling, endocrine pathways, and, eventually, metabolic and enzymatic changes.
Environmental
Environmental regulators of diapause generally display a characteristic seasonal pattern. In temperate regions, photoperiod is the most reliable cues of seasonal change. This informs entry into reproductive diapause for many northern insects, including the fruit fly Drosophila montana. Depending on the season in which diapause occurs, either short or long days can act as token stimuli. Insects may also respond to changing day length as well as relative day length. Temperature may also act as a regulating factor, either by inducing diapause or, more commonly, by modifying the response of the insect to photoperiod. Insects may respond to thermoperiod, the daily fluctuations of warm and cold that correspond with night and day, as well as to absolute or cumulative temperature. This has been observed in many moth species including the Indian mealmoth, where individuals diapause in different developmental stages due to environmental temperature. Food availability and quality may also help regulate diapause. In the desert locust, Schistocerca gregaria, a plant hormone called gibberellin stimulates reproductive development. During the dry season, when their food plants are in senescence and lacking gibberellin, the locusts remain immature and their reproductive tracts do not develop.
Neuroendocrine
The neuroendocrine system of insects consists primarily of neurosecretory cells in the brain, the corpora cardiaca, corpora allata and the prothoracic glands. There are several key hormones involved in the regulation of diapause: juvenile hormone (JH), diapause hormone (DH), and prothoracicotropic hormone (PTTH).
Prothoracicotropic hormone stimulates the prothoracic glands to produce ecdysteroids that are required to promote development. Larval and pupal diapauses are often regulated by an interruption of this connection, either by preventing release of prothoracicotropic hormone from the brain or by failure of the prothoracic glands to respond to prothoracicotropic hormone.
The corpora allata is responsible for the production of juvenile hormone (JH). In the bean bug, Riptortus pedestris, clusters of neurons on the protocerebrum called the pars lateralis maintain reproductive diapause by inhibiting JH production by the corpora allata. Adult diapause is often associated with the absence of JH, while larval diapause is often associated with its presence.
In adults, absence of JH causes degeneration of flight muscles and atrophy or cessation of development of reproductive tissues, and halts mating behaviour. The presence of JH in larvae may prevent moulting to the next larval instar, though successive stationary moults may still occur. In the corn borer, Diatraea gradiosella, JH is required for the accumulation by the fat body of a storage protein that is associated with diapause.
Diapause hormone regulates embryonic diapause in the eggs of the silkworm moth, Bombyx mori. DH is released from the subesophageal ganglion of the mother and triggers trehalase production by the ovaries. This generates high levels of glycogen in the eggs, which is converted into the polyhydric alcohols glycerol and sorbitol. Sorbitol directly inhibits the development of the embryos. Glycerol and sorbitol are reconverted into glycogen at the termination of diapause.
Tropical diapause
Diapause in the tropics is often initiated in response to biotic rather than abiotic components. For example, food in the form of vertebrate carcasses may be more abundant following dry seasons, or oviposition sites in the form of fallen trees may be more available following rainy seasons. Also, diapause may serve to synchronize mating seasons or reduce competition, rather than to avoid unfavourable climatic conditions.
Diapause in the tropics poses several challenges to insects that are not faced in temperate zones. Insects must reduce their metabolism without the aid of cold temperatures and may be faced with increased water loss due to high temperatures. While cold temperatures inhibit the growth of fungi and bacteria, diapausing tropical insects still have to deal with these pathogens. Also, predators and parasites may still be abundant during the diapause period.
Aggregations are common among diapausing tropical insects, especially in the orders Coleoptera, Lepidoptera, and Hemiptera. Aggregations may be used as protection against predation, since aggregating species are frequently toxic and predators quickly learn to avoid them. They can also serve to reduce water loss, as seen in the fungus beetle, Stenotarsus rotundus, which forms aggregations of up to 70,000 individuals, which may be eight beetles deep. Relative humidity is increased within the aggregations and beetles experience less water loss, probably due to decreased surface area to volume ratios reducing evaporative water loss.
| Biology and health sciences | Ethology | Biology |
1221466 | https://en.wikipedia.org/wiki/Coulometry | Coulometry | In analytical electrochemistry, coulometry is the measure of charge (coulombs) transfer during an electrochemical redox reaction. It can be used for precision measurements of charge, but coulometry is mainly used for analytical applications to determine the amount of matter transformed.
There are two main categories of coulometric techniques. Amperostatic coulometry, or coulometric titration keeps the current constant using an amperostat. Potentiostatic coulometry holds the electric potential constant during the reaction using a potentiostat.
History
The term coulometry was introduced in 1938 by Hungarian chemist László Szebellédy and Zoltan Somogyi. Coulometry is the measure of charge, thus named after its unit the coulomb.
Michael Faraday, known for his work in electricity and magnetism, made critical contributions to the field of electrochemistry. He discovered the laws of electrolysis, and in his recognition is the eponym of the Faraday constant. In the earliest developments of coulometry, Faraday proposed the first instrument to measure charge by utilizing the electrolysis of water.
Surface coulometry, the method of determining metallic layers or oxide films on metals, was first applied by American Chemist G. G. Grower in 1917 by checking the quality of tinned copper wire.
Coulometric methods were used widely in the middle of the twentieth century but voltammetric methods and non-electrochemical analytical methods took over decreasing the use for coulometry, but one method widely used today is the Karl Fischer method.
Potentiostatic coulometry
Potentiostatic coulometry utilizes a constant electric potential and is a technique most commonly referred to as "bulk electrolysis". Also called direct coulometry, the analyte is oxidized or reduced at the working electrode without intermediate reactions. The working electrode is kept at a constant potential and the current that flows through the circuit is measured. This constant potential is applied long enough to fully reduce or oxidize all of the electroactive species in a given solution. As the electroactive molecules are consumed, the current also decreases, approaching zero when the conversion is complete. The sample mass, molecular mass, number of electrons in the electrode reaction, and number of electrons passed during the experiment are all related by Faraday's laws. It follows that, if three of the values are known, then the fourth can be calculated.
Bulk electrolysis is often used to unambiguously assign the number of electrons consumed in a reaction observed through voltammetry. It also has the added benefit of producing a solution of a species (oxidation state) which may not be accessible through chemical routes. This species can then be isolated or further characterized while in solution.
The rate of such reactions is not determined by the concentration of the solution, but rather the mass transfer of the electroactive species in the solution to the electrode surface. Rates will increase when the volume of the solution is decreased, the solution is stirred more rapidly, or the area of the working electrode is increased. Since mass transfer is so important the solution is stirred during a bulk electrolysis. However, this technique is generally not considered a hydrodynamic technique, since a laminar flow of solution against the electrode is neither the objective nor outcome of the stirring.
The extent to which a reaction goes to completion is also related to how much greater the applied potential is than the reduction potential of interest. In the case where multiple reduction potentials are of interest, it is often difficult to set an electrolysis potential a "safe" distance (such as 200 mV) past a redox event. The result is incomplete conversion of the substrate, or else conversion of some of the substrate to the more reduced form. This factor must be considered when analyzing the current passed and when attempting to do further analysis/isolation/experiments with the substrate solution.
An advantage to this kind of analysis over electrogravimetry is that it does not require that the product of the reaction be weighed. This is useful for reactions where the product does not deposit as a solid, such as the determination of the amount of arsenic in a sample from the electrolysis of arsenous acid (H3AsO3) to arsenic acid (H3AsO4).
Coulometric titration
Coulometric titrations under a constant current system quantifies the to analyte by measuring the duration that current passes through the sample. In indirect or secondary coulometry, the working electrode produces a titrant that reacts with the analyte. When the analyte is completely consumed, endpoint detection is employed, preferably with an instrumental method for higher precision. The total charge that has flowed through the sample can be determined from the magnitude of the current (in amperes) and the duration of the current (in seconds). Using Faraday's Law, total charge can be used to determine the moles of the unknown species in solution. When the volume of the solution is known, the molarity of the unknown species can be determined.
Advantages of Coulometric Titration
Coulometric titration has the advantage that constant current sources for the generation of titrants are relatively easy to make.
The electrochemical generation of a titrant is much more sensitive and can be much more accurately controlled than the mechanical addition of titrant using a burette drive. For example, a constant current flow of 10 μA for 100 ms is easily generated and corresponds to about 10 micrograms of titrant.
The preparation of standard solutions and titer determination is no longer necessary.
Chemical substances that are unstable or difficult to handle because of their high volatility or reactivity in solution can also very easily be used as titrants. Examples are bromine, chlorine, Ti3+, Sn2+, Cr2+, and Karl Fischer reagents (iodine).
Coulometric titration can also be performed under inert atmosphere or be remotely controlled e.g. with radioactive substances.
Complete automation is simpler.
Applications
Karl Fischer reaction to determine water content
The Karl Fischer reaction uses a coulometric titration to determine the amount of water in a sample. It can determine concentrations of water on the order of milligrams per liter. It is used to find the amount of water in substances such as butter, sugar, cheese, paper, and petroleum.
The reaction involves converting solid iodine into hydrogen iodide in the presence of sulfur dioxide and water. Methanol is most often used as the solvent, but ethylene glycol and diethylene glycol also work. Pyridine is often used to prevent the buildup of sulfuric acid, although the use of imidazole and diethanolamine for this role are becoming more common. All reagents must be anhydrous for the analysis to be quantitative. The balanced chemical equation, using methanol and pyridine, is:
In this reaction, a single molecule of water reacts with a molecule of iodine. Since this technique is used to determine the water content of samples, atmospheric humidity could alter the results. Therefore, the system is usually isolated with drying tubes or placed in an inert gas container. In addition, the solvent will undoubtedly have some water in it so the solvent's water content must be measured to compensate for this inaccuracy.
To determine the amount of water in the sample, analysis must first be performed using either back or direct titration. In the direct method, just enough of the reagents will be added to completely use up all of the water. At this point in the titration, the current approaches zero. It is then possible to relate the amount of reagents used to the amount of water in the system via stoichiometry. The back-titration method is similar, but involves the addition of an excess of the reagent. This excess is then consumed by adding a known amount of a standard solution with known water content. The result reflects the water content of the sample and the standard solution. Since the amount of water in the standard solution is known, the difference reflects the water content of the sample.
Determination of film thickness
Coulometry can be used in the determination of the thickness of metallic coatings. This method is called surface coulometry and is performed by measuring the quantity of electricity needed to dissolve a well-defined area of the coating. The film thickness is proportional to the constant current , the molecular weight of the metal, the density of the metal, and the surface area :
The electrodes for this reaction are often platinum electrode and an electrode that relates to the reaction. For tin coating on a copper wire, a tin electrode is used, while a sodium chloride-zinc sulfate electrode would be used to determine the zinc film on a piece of steel. Special cells have been created to adhere to the surface of the metal to measure its thickness. These are basically columns with the internal electrodes with magnets or weights to attach to the surface. The results obtained by this coulometric method are similar to those achieved by other chemical and metallurgic techniques.
Coulometry in Healthcare
Determination of Chloride Levels
A type of clinical chemistry is measuring chloride levels in blood samples through a Cotlove chloridometer. Kidneys are responsible for the reabsorption of chloride to maintain electrolyte homeostasis. Measuring chloride levels allows for electrolyte stability, without this feature diseases such as hyperchoremia and hypochloremia would be harder to detect leaving body functions compromised.
Determination of Antioxidant Capacity in Human Blood
Coulometry can be used to measure the total antioxidant capacity (TAC) in blood and plasma through electrogenerated bromide. A method was developed that used TAC blood sampled from patients with chronic renal disease going through hemodialysis to research changes in TAC levels that could then be applied in clinics.
Coulometers
Electronic coulometer
The electronic coulometer is based on the application of the operational amplifier in the "integrator"-type circuit. The current passed through the resistor R1 makes a potential drop which is integrated by operational amplifier on the capacitor plates; the higher current, the larger the potential drop. The current need not be constant. In such scheme Vout is proportional of the passed charge. Sensitivity of the coulometer can be changed by choosing of the appropriate value of R1.
Electrochemical coulometers
There are three common types of coulometers based on electrochemical processes:
Copper coulometer
Mercury coulometer
Hofmann voltameter
"Voltameter" is a synonym for "coulometer".
Coulometric Microtitrators
An acid-base microtitorator utilizes the electrolysis of water, where protons or hydroxide ions are produced at the working electrode. The analyte reacts with the generated reagent, buffering the overall rate of reagent generation. A pH gradient forms from the diffusion of these reagents, where a pH sensor will determine the endpoint.
Some advantages of using a microtitrator include the fast completion time of the titration due to the micro-scale. Additionally, a negligibly small amount of the sample is consumed, so titrations can be repeatedly analyzed with the same sample. On the contrary, microtitrators require calibration because diffusion is variable, and thus this method is not absolute.
| Physical sciences | Electrical methods | Chemistry |
1222028 | https://en.wikipedia.org/wiki/Pine%20tar | Pine tar | Pine tar is a form of wood tar produced by the high temperature carbonization of pine wood in anoxic conditions (dry distillation or destructive distillation). The wood is rapidly decomposed by applying heat and pressure in a closed container; the primary resulting products are charcoal and pine tar.
Pine tar consists primarily of aromatic hydrocarbons, tar acids, and tar bases. Components of tar vary according to the pyrolytic process (e.g. method, duration, temperature) and origin of the wood (e.g. age of pine trees, type of soil, and moisture conditions during tree growth). The choice of wood, design of kiln, burning, and collection of the tar can vary. Only pine stumps and roots are used in the traditional production of pine tar.
Pine tar has a long history as a wood preservative, as a wood sealant for maritime use, in roofing construction and maintenance, in soaps, and in the treatment of carbuncles and skin diseases, such as psoriasis, eczema, and rosacea. It is used in baseball to enhance the grip of a hitter's bat; it is also sometimes used by pitchers to improve their grip on the ball, in violation of the rules.
History
Nordic Iron Age
Based on chemical analysis of organic residues, there are strong indications that cone-shaped pits discovered north of Uppsala, Sweden, have been used for pine tar production. Three of the pits have been radiocarbon dated. The oldest dates back to 540–380 BCE, which would make it the world's oldest known, still existing tar production facility. The other dates from 230–390 CE.
Classical antiquity
In his encyclopedic work Natural History () the Roman naturalist and author Pliny the Elder (23/24–79 CE) describes how, in Europe, tar is produced through the destructive distillation of pine wood. The wood was chopped into small pieces (billets) and heated in a furnace. The tar was used "for coating ships and for many other useful purposes."
17th century
Finland
Pine tar has long been used in Nordic nations as a preservative for wood which may be exposed to harsh conditions, including outdoor furniture and ship decking and rigging. Tar demand surged in the 17th century as European nations began constructing naval and merchant fleets that required tar and pitch for ship waterproofing and caulking. In Finland, then a part of Sweden, large-scale tar production began in the early 17th century. By the late 17th century, tar was Sweden's third most valuable export commodity.
North American colonies
Tar and pitch for maritime use was in such demand that it became an important export for the American colonies, which had extensive pine forests. North Carolinians became known as "Tar Heels."
18th–19th century
Sweden
In present-day Sweden, large-scale tar production started around 1700. Swedish pine tar was often called "Stockholm tar" since, for many years, a single company held a royal monopoly on its export out of Stockholm, Sweden. It was also known as "Archangel Tar". Stockholm tar became synonymous with top-quality tar. Tar production peaked around the mid-1850s and gradually declined thereafter, largely due to new shipbuilding materials reducing the need for tar as a waterproofing agent, and to the decreasing use of hemp rope as sailing ships were phased out.
Use
Pine tar was used as a preservative on the bottoms of traditional Nordic-style skis until modern synthetic materials replaced wood in their construction. It also helped waxes adhere, which aided such skis’ grip and glide.
Pine tar is widely used as a veterinary care product, particularly as an antiseptic and hoof care treatment for horses and cattle. It also has been used when chickens start pecking the low hen. Applying a smear of pine tar on the hens' wound acts as a natural germicidal/antibacterial agent that discourages continued attacks on the affected hen due to its foreign texture.
Pine tar is used as a softening solvent in the rubber industry, for treating and fabricating construction materials, and in special paints.
As a wood preservative
Pine tar is combined with gum turpentine and boiled linseed oil to create a wood preservative. First, a thin coat is applied using a mixture with a greater proportion of turpentine. This allows it to permeate deeper into the oakum and fibre of the wood and lets the tar seep into any pinholes and larger gaps that might be in the planks. The tar weeps out to the exterior and indicates where the boat needs the most attention. This is followed with a thicker standard mix. Such treatments, while effective, must be continually reapplied.
Weatherproofing rope
Traditionally, hemp and other natural fibers were the norm for rope production. Such rope would quickly rot when exposed to rain, and was typically tarred to preserve it. The tar would stain the hands of ship's crews, and British Navy seamen became known as "tars."
Baseball
Pine tar is applied to the handles of baseball bats to improve a batter's grip.
Rule 1.10(c) of the 2002 Official rules of Major League Baseball restricts application to the lower 18 inches of a bat. The most famous example of the rule being applied is the Pine Tar Incident, which occurred during the July 24, 1983 game between the Kansas City Royals and New York Yankees which resulted in a George Brett go-ahead home run in the ninth inning being nullified and the game being protested.
Pine tar is also sometimes used illegally by pitchers to improve their grip on the ball in cold weather. This is not allowed due to a regulation prohibiting the application of any foreign substance to a ball (except grip-improving baseball rubbing mud applied by the umpires).
Medical
Pine tar has historically been used for treating skin conditions, usually as an additive in cold process solid soap or lotions. Due to the high presence of phenol in the early manufacturing of pine tar, it was deemed carcinogenic. However, now much of the phenol has been removed. Pine tar was banned by the FDA along with many other ingredients categorized as over the counter drugs, due to a lack of evidence of safety and effectiveness for the specific uses named. However, clinical tests in Australia in 2017 demonstrated that the greatest risk comes from acute sensitivity for those with severe dermatological conditions, and if it comes in contact with the eyes. The number of positive reactions for wood tars was not significantly greater than those for other common allergens. In addition, the concentration of pine tar in topical products available in Australia is up to 2.3%, which is up to four times less than that tested in these studies.
Pine tar has been used to cover peck wounds in captive bird flocks such as chickens, to prevent continued pecking on a wound and cannibalism. Pine tar is also used in veterinary medicine as an expectorant and an antiseptic in chronic skin conditions.
| Physical sciences | Hydrocarbons | Chemistry |
6135482 | https://en.wikipedia.org/wiki/Juniper%20berry | Juniper berry | A juniper berry is the female seed cone produced by the various species of junipers. It is not a true berry but a cone with unusually fleshy and merged scales called a galbulus, which gives it a berry-like appearance. The cones from a handful of species, especially Juniperus communis, are used as a spice, particularly in European cuisine, and also give gin its distinctive flavour. Juniper berries are among the only spices derived from conifers, along with spruce buds.
Description
Unlike the separated and woody scales of a typical pine cone, those in a juniper berry remain fleshy and merge into a unified covering surrounding the seeds. Juniper berries are sometimes regarded as arils, like the berry-like cones of yews. Juniperus communis berries vary from to in diameter; other species are mostly similar in size, though some are larger, notably J. drupacea (). The berries are green when young and mature to purple-black over about 18 months in most species, including J. communis. Maturation occurs from as little as 8–10 months in some species up to over 24 months in J. drupacea. The mature, dark berries are usually (but not exclusively) used in cuisine, while gin is flavoured with fully grown, unripe berries.
Chemistry
Juniper berries contain diverse phytochemicals, including an essential oil in about 2% volume, a flavonoid called juniperin, resins (about 10% of volume), proteins, and acetic, malic and formic acids. From extracts of the berries, fatty acids, terpenes, aromatic compounds, and hydrocarbons, such as pinene, sabinene, terpinen-4-ol, limonene, and myrcene, were isolated.
Toxicity
While classified as generally recognized as safe in the United States, juniper berries may have various side effects that have not been tested extensively in clinical trials. Mainly due to an increased risk of miscarriage, even in small doses, consuming juniper berries may affect pregnant or breastfeeding women.
Allergic reactions are possible. Consuming large amounts of juniper berries may cause catharsis, convulsions, or harm kidney function. The berries of some species, such as J. sabina, are toxic.
Uses
The berries of some juniper species are considered too bitter to eat. In addition to J. communis and J. drupacea, edible species include J. phoenicea, J. deppeana, and J. californica.
The flavour profile of young, green berries is dominated by pinene; as they mature this piney, resinous backdrop is joined by what Harold McGee describes as "green-fresh" and citrus notes. The outer scales of the berries are relatively flavourless, so the berries are almost always at least lightly crushed before being used as a spice. They are used both fresh and dried, but their flavour and odour are at their strongest immediately after harvest and decline during drying and storage.
Flavour
Juniper berries are used in northern European and particularly Scandinavian cuisine to, according to one source, "impart a sharp, clear flavor" to meat dishes, especially wild birds (including thrush, blackbird, and woodcock) and game meats (including boar and venison). They also season pork, cabbage, and sauerkraut dishes. Traditional recipes for choucroute garnie, an Alsatian dish of sauerkraut and meats, universally include juniper berries. Besides Norwegian, Danish and Swedish dishes, juniper berries are also sometimes used in German, Austrian, Czech, Polish and Hungarian cuisine, often with roasts (such as German ). Northern Italian cuisine, especially that of the South Tyrol, also incorporates juniper berries. They are also used in the Italian region of Apulia, especially to flavour brines.
Juniper, typically J. communis, is used to flavor gin, a liquor developed in the 17th century in the Netherlands. The name gin itself is derived from either the French or the Dutch , both of which mean "juniper". Other juniper-flavoured beverages include the Finnish rye-and-juniper beer known as sahti, which is flavored with both juniper berries and branches.
Another drink made from the berries is a , a soft drink made in Sweden mainly sold during Christmas.
Food
A few North American juniper species produce a seed cone with a sweeter, less resinous flavor than those typically used as a spice. For example, one field guide describes the flesh of the berries of J. californica as "dry, mealy, and fibrous but sweet and without resin cells". Such species have been used not just as a seasoning but as a nutritive food by some Native Americans. The berries also have medicinal uses. For example, the Blackfoot used juniper berry tea to cure vomiting, while Crow women drank juniper berry tea after childbirth to increase cleansing and healing. In addition to medicinal and culinary purposes, Native Americans have also used the seeds inside juniper berries as beads for jewellery and decoration.
An essential oil extracted from juniper berries is used in aromatherapy, both for body massage, diffusion, and perfumery.
Culture
Juniper berries, including Juniperus phoenicea and J. oxycedrus, have been found in ancient Egyptian tombs at multiple sites. J. oxycedrus is not known to grow in Egypt, and neither is J. excelsa, which was found along with J. oxycedrus in the tomb of Tutankhamun. The berries imported into Egypt may have come from Greece; the Greeks record using juniper berries as a medicine long before mentioning their use in food.
The Greeks used the berries in many of their Olympics events because of their belief that the berries increased physical stamina in athletes.
The Romans used juniper berries as a cheap domestically produced substitute for the expensive black pepper and long pepper imported from India. It was also used as an adulterant, as reported in Pliny the Elder's Natural History: "pepper is adulterated with juniper berries, which have the property, to a marvellous degree, of assuming the pungency of pepper". Pliny also incorrectly asserted that black pepper grew on trees that were "very similar in appearance to our junipers".
The berries were an integral part of Desert Serrano (Vanyume) culture and grew throughout the Mojave River region. The major village of Wá’peat was derived from the Serrano word for juniper berries, .
| Biology and health sciences | Cupressaceae | Plants |
6139438 | https://en.wikipedia.org/wiki/Formation%20and%20evolution%20of%20the%20Solar%20System | Formation and evolution of the Solar System | There is evidence that the formation of the Solar System began about 4.6 billion years ago with the gravitational collapse of a small part of a giant molecular cloud. Most of the collapsing mass collected in the center, forming the Sun, while the rest flattened into a protoplanetary disk out of which the planets, moons, asteroids, and other small Solar System bodies formed.
This model, known as the nebular hypothesis, was first developed in the 18th century by Emanuel Swedenborg, Immanuel Kant, and Pierre-Simon Laplace. Its subsequent development has interwoven a variety of scientific disciplines including astronomy, chemistry, geology, physics, and planetary science. Since the dawn of the Space Age in the 1950s and the discovery of exoplanets in the 1990s, the model has been both challenged and refined to account for new observations.
The Solar System has evolved considerably since its initial formation. Many moons have formed from circling discs of gas and dust around their parent planets, while other moons are thought to have formed independently and later to have been captured by their planets. Still others, such as Earth's Moon, may be the result of giant collisions. Collisions between bodies have occurred continually up to the present day and have been central to the evolution of the Solar System. Beyond Neptune, many sub-planet sized objects formed. Several thousand trans-Neptunian objects have been observed. Unlike the planets, these trans-Neptunian objects mostly move on eccentric orbits, inclined to the plane of the planets. The positions of the planets might have shifted due to gravitational interactions. Planetary migration may have been responsible for much of the Solar System's early evolution.
In roughly 5 billion years, the Sun will cool and expand outward to many times its current diameter (becoming a red giant), before casting off its outer layers as a planetary nebula and leaving behind a stellar remnant known as a white dwarf. In the distant future, the gravity of passing stars will gradually reduce the Sun's retinue of planets. Some planets will be destroyed, and others ejected into interstellar space. Ultimately, over the course of tens of billions of years, it is likely that the Sun will be left with none of the original bodies in orbit around it.
History
Ideas concerning the origin and fate of the world date from the earliest known writings; however, for almost all of that time, there was no attempt to link such theories to the existence of a "Solar System", simply because it was not generally thought that the Solar System, in the sense we now understand it, existed. The first step toward a theory of Solar System formation and evolution was the general acceptance of heliocentrism, which placed the Sun at the centre of the system and the Earth in orbit around it. This concept had been developed for millennia (Aristarchus of Samos had suggested it as early as 250 BC), but was not widely accepted until the end of the 17th century. The first recorded use of the term "Solar System" dates from 1704.
The current standard theory for Solar System formation, the nebular hypothesis, has fallen into and out of favour since its formulation by Emanuel Swedenborg, Immanuel Kant, and Pierre-Simon Laplace in the 18th century. The most significant criticism of the hypothesis was its apparent inability to explain the Sun's relative lack of angular momentum when compared to the planets. However, since the early 1980s studies of young stars have shown them to be surrounded by cool discs of dust and gas, exactly as the nebular hypothesis predicts, which has led to its re-acceptance.
Understanding of how the Sun is expected to continue to evolve required an understanding of the source of its power. Arthur Stanley Eddington's confirmation of Albert Einstein's theory of relativity led to his realisation that the Sun's energy comes from nuclear fusion reactions in its core, fusing hydrogen into helium. In 1935, Eddington went further and suggested that other elements also might form within stars. Fred Hoyle elaborated on this premise by arguing that evolved stars called red giants created many elements heavier than hydrogen and helium in their cores. When a red giant finally casts off its outer layers, these elements would then be recycled to form other star systems.
Formation
Presolar nebula
The nebular hypothesis says that the Solar System formed from the gravitational collapse of a fragment of a giant molecular cloud, most likely at the edge of a Wolf-Rayet bubble. The cloud was about 20 parsecs (65 light years) across, while the fragments were roughly 1 parsec (three and a quarter light-years) across. The further collapse of the fragments led to the formation of dense cores 0.01–0.1 parsec (2,000–20,000 AU) in size. One of these collapsing fragments (known as the presolar nebula) formed what became the Solar System. The composition of this region with a mass just over that of the Sun () was about the same as that of the Sun today, with hydrogen, along with helium and trace amounts of lithium produced by Big Bang nucleosynthesis, forming about 98% of its mass. The remaining 2% of the mass consisted of heavier elements that were created by nucleosynthesis in earlier generations of stars. Late in the life of these stars, they ejected heavier elements into the interstellar medium. Some scientists have given the name Coatlicue to a hypothetical star that went supernova and created the presolar nebula.
The oldest inclusions found in meteorites, thought to trace the first solid material to form in the presolar nebula, are 4,568.2 million years old, which is one definition of the age of the Solar System. Studies of ancient meteorites reveal traces of stable daughter nuclei of short-lived isotopes, such as iron-60, that only form in exploding, short-lived stars. This indicates that one or more supernovae occurred nearby. A shock wave from a supernova may have triggered the formation of the Sun by creating relatively dense regions within the cloud, causing these regions to collapse. The highly homogeneous distribution of iron-60 in the Solar System points to the occurrence of this supernova and its injection of iron-60 being well before the accretion of nebular dust into planetary bodies. Because only massive, short-lived stars produce supernovae, the Sun must have formed in a large star-forming region that produced massive stars, possibly similar to the Orion Nebula. Studies of the structure of the Kuiper belt and of anomalous materials within it suggest that the Sun formed within a cluster of between 1,000 and 10,000 stars with a diameter of between 6.5 and 19.5 light years and a collective mass of . This cluster began to break apart between 135 million and 535 million years after formation. Several simulations of our young Sun interacting with close-passing stars over the first 100 million years of its life produced anomalous orbits observed in the outer Solar System, such as detached objects. A recent study suggests that such a passing star is not only responsible for the orbits of the detached objects but also the hot and cold Kuiper belt population, the Sedna-like objects, the extreme TNOs and the retrograde TNOs.
Because of the conservation of angular momentum, the nebula spun faster as it collapsed. As the material within the nebula condensed, the temperature rose. The center, where most of the mass collected, became increasingly hotter than the surrounding disc. Over about 100,000 years, the competing forces of gravity, gas pressure, magnetic fields, and rotation caused the contracting nebula to flatten into a spinning protoplanetary disc with a diameter of about 200 AU and form a hot, dense protostar (a star in which hydrogen fusion has not yet begun) at the centre. Since about half of all known stars form systems of multiple stars and because Jupiter is made of the same elements as the Sun (hydrogen and helium), it has been suggested that the Solar System might have been early in its formation a protostar system with Jupiter being the second but failed protostar, but Jupiter has far too little mass to trigger fusion in its core and so became a gas giant; it is in fact younger than the Sun and the oldest planet of the Solar System.
At this point in the Sun's evolution, the Sun is thought to have been a T Tauri star. Studies of T Tauri stars show that they are often accompanied by discs of pre-planetary matter with masses of . These discs extend to several hundred AU—the Hubble Space Telescope has observed protoplanetary discs of up to 1000 AU in diameter in star-forming regions such as the Orion Nebula—and are rather cool, reaching a surface temperature of only about at their hottest.
Within 50 million years, the temperature and pressure at the core of the Sun became so great that its hydrogen began to fuse, creating an internal source of energy that countered gravitational contraction until hydrostatic equilibrium was achieved. This marked the Sun's entry into the prime phase of its life, known as the main sequence. Main-sequence stars derive energy from the fusion of hydrogen into helium in their cores. The Sun remains a main-sequence star today.
As the early Solar System continued to evolve, it eventually drifted away from its siblings in the stellar nursery, and continued orbiting the Milky Way's center on its own. The Sun likely drifted from its original orbital distance from the center of the galaxy. The chemical history of the Sun suggests it may have formed as much as 3 kpc closer to the galaxy core.
Solar System birth environment
Like most stars, the Sun likely formed not in isolation but as part of a young star cluster. There are several indications that hint at the cluster environment having had some influence over the young, still-forming Solar System. For example, the decline in mass beyond Neptune and the extreme eccentric-orbit of Sedna have been interpreted as a signature of the Solar System having been influenced by its birth environment. Whether the presence of the isotopes iron-60 and aluminium-26 can be interpreted as a sign of a birth cluster containing massive stars is still under debate. If the Sun was part of a star cluster, it might have been influenced by close flybys of other stars, the strong radiation of nearby massive stars and ejecta from supernovae occurring close by.
Formation of the planets
The various planets are thought to have formed from the solar nebula, the disc-shaped cloud of gas and dust left over from the Sun's formation. The currently accepted method by which the planets formed is accretion, in which the planets began as dust grains in orbit around the central protostar. Through direct contact and self-organization, these grains formed into clumps up to in diameter, which in turn collided to form larger bodies (planetesimals) of ~ in size. These gradually increased through further collisions, growing at the rate of centimetres per year over the course of the next few million years.
The inner Solar System, the region of the Solar System inside 4 AU, was too warm for volatile molecules like water and methane to condense, so the planetesimals that formed there could only form from compounds with high melting points, such as metals (like iron, nickel, and aluminium) and rocky silicates. These rocky bodies would become the terrestrial planets (Mercury, Venus, Earth, and Mars). These compounds are quite rare in the Universe, comprising only 0.6% of the mass of the nebula, so the terrestrial planets could not grow very large. The terrestrial embryos grew to about 0.05 Earth masses () and ceased accumulating matter about 100,000 years after the formation of the Sun; subsequent collisions and mergers between these planet-sized bodies allowed terrestrial planets to grow to their present sizes.
When terrestrial planets were forming, they remained immersed in a disk of gas and dust. Pressure partially supported the gas and so did not orbit the Sun as rapidly as the planets. The resulting drag and, more importantly, gravitational interactions with the surrounding material caused a transfer of angular momentum, and as a result the planets gradually migrated to new orbits. Models show that density and temperature variations in the disk governed this rate of migration, but the net trend was for the inner planets to migrate inward as the disk dissipated, leaving the planets in their current orbits.
The giant planets (Jupiter, Saturn, Uranus, and Neptune) formed further out, beyond the frost line, which is the point between the orbits of Mars and Jupiter where the material is cool enough for volatile icy compounds to remain solid. The ices that formed the Jovian planets were more abundant than the metals and silicates that formed the terrestrial planets, allowing the giant planets to grow massive enough to capture hydrogen and helium, the lightest and most abundant elements. Planetesimals beyond the frost line accumulated up to within about 3 million years. Today, the four giant planets comprise just under 99% of all the mass orbiting the Sun. Theorists believe it is no accident that Jupiter lies just beyond the frost line. Because the frost line accumulated large amounts of water via evaporation from infalling icy material, it created a region of lower pressure that increased the speed of orbiting dust particles and halted their motion toward the Sun. In effect, the frost line acted as a barrier that caused the material to accumulate rapidly at ~5 AU from the Sun. This excess material coalesced into a large embryo (or core) on the order of , which began to accumulate an envelope via accretion of gas from the surrounding disc at an ever-increasing rate. Once the envelope mass became about equal to the solid core mass, growth proceeded very rapidly, reaching about 150 Earth masses ~105 years thereafter and finally topping out at . Saturn may owe its substantially lower mass simply to having formed a few million years after Jupiter, when there was less gas available to consume.
T Tauri stars like the young Sun have far stronger stellar winds than more stable, older stars. Uranus and Neptune are thought to have formed after Jupiter and Saturn did, when the strong solar wind had blown away much of the disc material. As a result, those planets accumulated little hydrogen and helium—not more than each. Uranus and Neptune are sometimes referred to as failed cores. The main problem with formation theories for these planets is the timescale of their formation. At the current locations it would have taken millions of years for their cores to accrete. This means that Uranus and Neptune may have formed closer to the Sun—near or even between Jupiter and Saturn—later migrating or being ejected outward (see Planetary migration below). Motion in the planetesimal era was not all inward toward the Sun; the Stardust sample return from Comet Wild 2 has suggested that materials from the early formation of the Solar System migrated from the warmer inner Solar System to the region of the Kuiper belt.
After between three and ten million years, the young Sun's solar wind would have cleared away all the gas and dust in the protoplanetary disc, blowing it into interstellar space, thus ending the growth of the planets.
Subsequent evolution
The planets were originally thought to have formed in or near their current orbits. This has been questioned during the last 20 years. Currently, many planetary scientists think that the Solar System might have looked very different after its initial formation: several objects at least as massive as Mercury may have been present in the inner Solar System, the outer Solar System may have been much more compact than it is now, and the Kuiper belt may have been much closer to the Sun.
Terrestrial planets
At the end of the planetary formation epoch, the inner Solar System was populated by 50–100 Moon-to-Mars-sized protoplanets. Further growth was possible only because these bodies collided and merged, which took less than 100 million years. These objects would have gravitationally interacted with one another, tugging at each other's orbits until they collided, growing larger until the four terrestrial planets we know today took shape. One such giant collision is thought to have formed the Moon (see Moons below), while another removed the outer envelope of the young Mercury.
One unresolved issue with this model is that it cannot explain how the initial orbits of the proto-terrestrial planets, which would have needed to be highly eccentric in order to collide, produced the remarkably stable and nearly circular orbits they have today. One hypothesis for this "eccentricity dumping" is that terrestrials formed in a disc of gas still not expelled by the Sun. The "gravitational drag" of this residual gas would have eventually lowered the planets' energy, smoothing out their orbits. However, such gas, if it existed, would have prevented the terrestrial planets' orbits from becoming so eccentric in the first place. Another hypothesis is that gravitational drag occurred not between the planets and residual gas but between the planets and the remaining small bodies. As the large bodies moved through the crowd of smaller objects, the smaller objects, attracted by the larger planets' gravity, formed a region of higher density, a "gravitational wake", in the larger objects' path. As they did so, the increased gravity of the wake slowed the larger objects down into more regular orbits.
Asteroid belt
The outer edge of the terrestrial region, between 2 and 4 AU from the Sun, is called the asteroid belt. The asteroid belt initially contained more than enough matter to form 2–3 Earth-like planets, and, indeed, a large number of planetesimals formed there. As with the terrestrials, planetesimals in this region later coalesced and formed 20–30 Moon- to Mars-sized planetary embryos; however, the proximity of Jupiter meant that after this planet formed, 3 million years after the Sun, the region's history changed dramatically. Orbital resonances with Jupiter and Saturn are particularly strong in the asteroid belt, and gravitational interactions with more massive embryos scattered many planetesimals into those resonances. Jupiter's gravity increased the velocity of objects within these resonances, causing them to shatter upon collision with other bodies, rather than accrete.
As Jupiter migrated inward following its formation (see Planetary migration below), resonances would have swept across the asteroid belt, dynamically exciting the region's population and increasing their velocities relative to each other. The cumulative action of the resonances and the embryos either scattered the planetesimals away from the asteroid belt or excited their orbital inclinations and eccentricities. Some of those massive embryos too were ejected by Jupiter, while others may have migrated to the inner Solar System and played a role in the final accretion of the terrestrial planets. During this primary depletion period, the effects of the giant planets and planetary embryos left the asteroid belt with a total mass equivalent to less than 1% that of the Earth, composed mainly of small planetesimals.
This is still 10–20 times more than the current mass in the main belt, which is now about . A secondary depletion period that brought the asteroid belt down close to its present mass is thought to have followed when Jupiter and Saturn entered a temporary 2:1 orbital resonance (see below).
The inner Solar System's period of giant impacts probably played a role in Earth acquiring its current water content (~6 kg) from the early asteroid belt. Water is too volatile to have been present at Earth's formation and must have been subsequently delivered from outer, colder parts of the Solar System. The water was probably delivered by planetary embryos and small planetesimals thrown out of the asteroid belt by Jupiter. A population of main-belt comets discovered in 2006 has also been suggested as a possible source for Earth's water. In contrast, comets from the Kuiper belt or farther regions delivered not more than about 6% of Earth's water. The panspermia hypothesis holds that life itself may have been deposited on Earth in this way, although this idea is not widely accepted.
Planetary migration
According to the nebular hypothesis, the outer two planets may be in the "wrong place". Uranus and Neptune (known as the "ice giants") exist in a region where the reduced density of the solar nebula and longer orbital times render their formation there highly implausible. The two are instead thought to have formed in orbits near Jupiter and Saturn (known as the "gas giants"), where more material was available, and to have migrated outward to their current positions over hundreds of millions of years.
The migration of the outer planets is also necessary to account for the existence and properties of the Solar System's outermost regions. Beyond Neptune, the Solar System continues into the Kuiper belt, the scattered disc, and the Oort cloud, three sparse populations of small icy bodies thought to be the points of origin for most observed comets. At their distance from the Sun, accretion was too slow to allow planets to form before the solar nebula dispersed, and thus the initial disc lacked enough mass density to consolidate into a planet. The Kuiper belt lies between 30 and 55 AU from the Sun, while the farther scattered disc extends to over 100 AU, and the distant Oort cloud begins at about 50,000 AU. Originally, however, the Kuiper belt was much denser and closer to the Sun, with an outer edge at approximately 30 AU. Its inner edge would have been just beyond the orbits of Uranus and Neptune, which were in turn far closer to the Sun when they formed (most likely in the range of 15–20 AU), and in 50% of simulations ended up in opposite locations, with Uranus farther from the Sun than Neptune.
According to the Nice model, after the formation of the Solar System, the orbits of all the giant planets continued to change slowly, influenced by their interaction with the large number of remaining planetesimals. After 500–600 million years (about 4 billion years ago) Jupiter and Saturn fell into a 2:1 resonance: Saturn orbited the Sun once for every two Jupiter orbits. This resonance created a gravitational push against the outer planets, possibly causing Neptune to surge past Uranus and plough into the ancient Kuiper belt.
The planets scattered the majority of the small icy bodies inwards, while themselves moving outwards. These planetesimals then scattered off the next planet they encountered in a similar manner, moving the planets' orbits outwards while they moved inwards. This process continued until the planetesimals interacted with Jupiter, whose immense gravity sent them into highly elliptical orbits or even ejected them outright from the Solar System. This caused Jupiter to move slightly inward. Those objects scattered by Jupiter into highly elliptical orbits formed the Oort cloud; those objects scattered to a lesser degree by the migrating Neptune formed the current Kuiper belt and scattered disc. This scenario explains the Kuiper belt's and scattered disc's present low mass. Some of the scattered objects, including Pluto, became gravitationally tied to Neptune's orbit, forcing them into mean-motion resonances. Eventually, friction within the planetesimal disc made the orbits of Uranus and Neptune near-circular again.
In contrast to the outer planets, the inner planets are not thought to have migrated significantly over the age of the Solar System, because their orbits have remained stable following the period of giant impacts.
Another question is why Mars came out so small compared with Earth. A study by Southwest Research Institute, San Antonio, Texas, published June 6, 2011 (called the Grand tack hypothesis), proposes that Jupiter had migrated inward to 1.5 AU. After Saturn formed, migrated inward, and established the 2:3 mean motion resonance with Jupiter, the study assumes that both planets migrated back to their present positions. Jupiter thus would have consumed much of the material that would have created a bigger Mars. The same simulations also reproduce the characteristics of the modern asteroid belt, with dry asteroids and water-rich objects similar to comets. However, it is unclear whether conditions in the solar nebula would have allowed Jupiter and Saturn to move back to their current positions, and according to current estimates this possibility appears unlikely. Moreover, alternative explanations for the small mass of Mars exist.
Late Heavy Bombardment and after
Gravitational disruption from the outer planets' migration would have sent large numbers of asteroids into the inner Solar System, severely depleting the original belt until it reached today's extremely low mass. This event may have triggered the Late Heavy Bombardment that is hypothesised to have occurred approximately 4 billion years ago, 500–600 million years after the formation of the Solar System. However, a recent re-appraisal of the cosmo-chemical constraints indicates that there was likely no late spike (“terminal cataclysm”) in the bombardment rate.
If it occurred, this period of heavy bombardment lasted several hundred million years and is evident in the cratering still visible on geologically dead bodies of the inner Solar System such as the Moon and Mercury. The oldest known evidence for life on Earth dates to 3.8 billion years ago—almost immediately after the end of the Late Heavy Bombardment.
Impacts are thought to be a regular (if currently infrequent) part of the evolution of the Solar System. That they continue to happen is evidenced by the collision of Comet Shoemaker–Levy 9 with Jupiter in 1994, the 2009 Jupiter impact event, the Tunguska event, the Chelyabinsk meteor and the impact that created Meteor Crater in Arizona. The process of accretion, therefore, is not complete, and may still pose a threat to life on Earth.
Over the course of the Solar System's evolution, comets were ejected out of the inner Solar System by the gravity of the giant planets and sent thousands of AU outward to form the Oort cloud, a spherical outer swarm of cometary nuclei at the farthest extent of the Sun's gravitational pull. Eventually, after about 800 million years, the gravitational disruption caused by galactic tides, passing stars and giant molecular clouds began to deplete the cloud, sending comets into the inner Solar System. The evolution of the outer Solar System also appears to have been influenced by space weathering from the solar wind, micrometeorites, and the neutral components of the interstellar medium.
The evolution of the asteroid belt after Late Heavy Bombardment was mainly governed by collisions. Objects with large mass have enough gravity to retain any material ejected by a violent collision. In the asteroid belt this usually is not the case. As a result, many larger objects have been broken apart, and sometimes newer objects have been forged from the remnants in less violent collisions. Moons around some asteroids currently can only be explained as consolidations of material flung away from the parent object without enough energy to entirely escape its gravity.
Moons
Moons have come to exist around most planets and many other Solar System bodies. These natural satellites originated by one of three possible mechanisms:
Co-formation from a circumplanetary disc (only in the cases of the giant planets);
Formation from impact debris (given a large enough impact at a shallow angle); and
Capture of a passing object.
Jupiter and Saturn have several large moons, such as Io, Europa, Ganymede and Titan, which may have originated from discs around each giant planet in much the same way that the planets formed from the disc around the Sun. This origin is indicated by the large sizes of the moons and their proximity to the planet. These attributes are impossible to achieve via capture, while the gaseous nature of the primaries also makes formation from collision debris unlikely. The outer moons of the giant planets tend to be small and have eccentric orbits with arbitrary inclinations. These are the characteristics expected of captured bodies. Most such moons orbit in the direction opposite to the rotation of their primary. The largest irregular moon is Neptune's moon Triton, which is thought to be a captured Kuiper belt object.
Moons of solid Solar System bodies have been created by both collisions and capture. Mars's two small moons, Deimos and Phobos, are thought to be captured asteroids.
The Earth's moon is thought to have formed as a result of a single, large head-on collision.
The impacting object probably had a mass comparable to that of Mars, and the impact probably occurred near the end of the period of giant impacts. The collision kicked into orbit some of the impactor's mantle, which then coalesced into the Moon. The impact was probably the last in a series of mergers that formed the Earth.
It has been further hypothesized that the Mars-sized object may have formed at one of the stable Earth–Sun Lagrangian points (either or ) and drifted from its position. The moons of trans-Neptunian objects Pluto (Charon) and Orcus (Vanth) may also have formed by means of a large collision: the Pluto–Charon, Orcus–Vanth and Earth–Moon systems are unusual in the Solar System in that the satellite's mass is at least 1% that of the larger body.
Future
Astronomers estimate that the current state of the Solar System will not change drastically until the Sun has fused almost all the hydrogen fuel in its core into helium, beginning its evolution from the main sequence of the Hertzsprung–Russell diagram and into its red-giant phase. The Solar System will continue to evolve until then. Eventually, the Sun will likely expand sufficiently to overwhelm the inner planets (Mercury, Venus, and possibly Earth) but not the outer planets, including Jupiter and Saturn. Afterward, the Sun would be reduced to the size of a white dwarf, and the outer planets and their moons would continue orbiting this diminutive solar remnant. This future development may be similar to the observed detection of MOA-2010-BLG-477L b, a Jupiter-sized exoplanet orbiting its host white dwarf star MOA-2010-BLG-477L.
Long-term stability
The Solar System is chaotic over million- and billion-year timescales, with the orbits of the planets open to long-term variations. One notable example of this chaos is the Neptune–Pluto system, which lies in a 3:2 orbital resonance. Although the resonance itself will remain stable, it becomes impossible to predict the position of Pluto with any degree of accuracy more than 10–20 million years (the Lyapunov time) into the future. Another example is Earth's axial tilt, which, due to friction raised within Earth's mantle by tidal interactions with the Moon (see below), is incomputable from some point between 1.5 and 4.5 billion years from now.
The outer planets' orbits are chaotic over longer timescales, with a Lyapunov time in the range of 2–230 million years.
In all cases, this means that the position of a planet along its orbit ultimately becomes impossible to predict with any certainty (so, for example, the timing of winter and summer becomes uncertain). Still, in some cases, the orbits themselves may change dramatically. Such chaos manifests most strongly as changes in eccentricity, with some planets' orbits becoming significantly more—or less—elliptical.
Ultimately, the Solar System is stable in that none of the planets are likely to collide with each other or be ejected from the system in the next few billion years. Beyond this, within five billion years or so, Mars's eccentricity may grow to around 0.2, such that it lies on an Earth-crossing orbit, leading to a potential collision. In the same timescale, Mercury's eccentricity may grow even further, and a close encounter with Venus could theoretically eject it from the Solar System altogether or send it on a collision course with Venus or Earth. This could happen within a billion years, according to numerical simulations in which Mercury's orbit is perturbed.
Moon–ring systems
The evolution of moon systems is driven by tidal forces. A moon will raise a tidal bulge in the object it orbits (the primary) due to the differential gravitational force across diameter of the primary. If a moon is revolving in the same direction as the planet's rotation and the planet is rotating faster than the orbital period of the moon, the bulge will constantly be pulled ahead of the moon. In this situation, angular momentum is transferred from the rotation of the primary to the revolution of the satellite. The moon gains energy and gradually spirals outward, while the primary rotates more slowly over time.
The Earth and its Moon are one example of this configuration. Today, the Moon is tidally locked to the Earth; one of its revolutions around the Earth (currently about 29 days) is equal to one of its rotations about its axis, so it always shows one face to the Earth. The Moon will continue to recede from Earth, and Earth's spin will continue to slow gradually. Other examples are the Galilean moons of Jupiter (as well as many of Jupiter's smaller moons) and most of the larger moons of Saturn.
A different scenario occurs when the moon is either revolving around the primary faster than the primary rotates or is revolving in the direction opposite the planet's rotation. In these cases, the tidal bulge lags behind the moon in its orbit. In the former case, the direction of angular momentum transfer is reversed, so the rotation of the primary speeds up while the satellite's orbit shrinks. In the latter case, the angular momentum of the rotation and revolution have opposite signs, so transfer leads to decreases in the magnitude of each (that cancel each other out). In both cases, tidal deceleration causes the moon to spiral in towards the primary until it either is torn apart by tidal stresses, potentially creating a planetary ring system, or crashes into the planet's surface or atmosphere. Such a fate awaits the moons Phobos of Mars (within 30 to 50 million years), Triton of Neptune (in 3.6 billion years), and at least 16 small satellites of Uranus and Neptune. Uranus's Desdemona may even collide with one of its neighboring moons.
A third possibility is where the primary and moon are tidally locked to each other. In that case, the tidal bulge stays directly under the moon, there is no angular momentum transfer, and the orbital period will not change. Pluto and Charon are an example of this type of configuration.
There is no consensus on the mechanism of the formation of the rings of Saturn. Although theoretical models indicated that the rings were likely to have formed early in the Solar System's history, data from the Cassini–Huygens spacecraft suggests they formed relatively late.
The Sun and planetary environments
In the long term, the greatest changes in the Solar System will come from changes in the Sun itself as it ages. As the Sun burns through its hydrogen fuel supply, it gets hotter and burns the remaining fuel even faster. As a result, the Sun is growing brighter at a rate of ten percent every 1.1 billion years. In about 600 million years, the Sun's brightness will have disrupted the Earth's carbon cycle to the point where trees and forests (C3 photosynthetic plant life) will no longer be able to survive; and in around 800 million years, the Sun will have killed all complex life on the Earth's surface and in the oceans. In 1.1 billion years, the Sun's increased radiation output will cause its circumstellar habitable zone to move outwards, making the Earth's surface too hot for liquid water to exist there naturally. At this point, all life will be reduced to single-celled organisms. Evaporation of water, a potent greenhouse gas, from the oceans' surface could accelerate temperature increase, potentially ending all life on Earth even sooner. During this time, it is possible that as Mars's surface temperature gradually rises, carbon dioxide and water currently frozen under the surface regolith will release into the atmosphere, creating a greenhouse effect that will heat the planet until it achieves conditions parallel to Earth today, providing a potential future abode for life. By 3.5 billion years from now, Earth's surface conditions will be similar to those of Venus today.
Around 5.4 billion years from now, the core of the Sun will become hot enough to trigger hydrogen fusion in its surrounding shell. This will cause the outer layers of the star to expand greatly, and the star will enter a phase of its life in which it is called a red giant. Within 7.5 billion years, the Sun will have expanded to a radius of —256 times its current size. At the tip of the red-giant branch, as a result of the vastly increased surface area, the Sun's surface will be much cooler (about ) than now, and its luminosity much higher—up to 2,700 current solar luminosities. For part of its red-giant life, the Sun will have a strong stellar wind that will carry away around 33% of its mass. During these times, it is possible that Saturn's moon Titan could achieve surface temperatures necessary to support life.
As the Sun expands, it will swallow the planets Mercury and Venus. Earth's fate is less clear; although the Sun will envelop Earth's current orbit, the star's loss of mass (and thus weaker gravity) will cause the planets' orbits to move farther out. If it were only for this, Venus and Earth would probably escape incineration, but a 2008 study suggests that Earth will likely be swallowed up as a result of tidal interactions with the Sun's weakly-bound outer envelope.
Additionally, the Sun's habitable zone will move into the outer Solar System and eventually beyond the Kuiper belt at the end of the red-giant phase, causing icy bodies such as Enceladus and Pluto to thaw. During this time, these worlds could support a water-based hydrologic cycle, but as they were too small to hold a dense atmosphere like Earth, they would experience extreme day–night temperature differences. When the Sun leaves the red-giant branch and enters the asymptotic giant branch, the habitable zone will abruptly shrink to roughly the space between Jupiter and Saturn's present-day orbits, but toward the end of the 200 million-year duration of the asymptotic giant phase, it will expand outward to about the same distance as before.
Gradually, the hydrogen burning in the shell around the solar core will increase the mass of the core until it reaches about 45% of the present solar mass. At this point, the density and temperature will become so high that the fusion of helium into carbon will begin, leading to a helium flash; the Sun will shrink from around 250 to 11 times its present (main-sequence) radius. Consequently, its luminosity will decrease from around 3,000 to 54 times its current level, and its surface temperature will increase to about . The Sun will become a horizontal giant, burning helium in its core in a stable fashion, much like it burns hydrogen today. The helium-fusing stage will last only 100 million years. Eventually, it will have to again resort to the reserves of hydrogen and helium in its outer layers. It will expand a second time, becoming what is known as an asymptotic giant. Here the luminosity of the Sun will increase again, reaching about 2,090 present luminosities, and it will cool to about . This phase lasts about 30 million years, after which, over the course of a further 100,000 years, the Sun's remaining outer layers will fall away, ejecting a vast stream of matter into space and forming a halo known (misleadingly) as a planetary nebula. The ejected material will contain the helium and carbon produced by the Sun's nuclear reactions, continuing the enrichment of the interstellar medium with heavy elements for future generations of stars and planets.
This is a relatively peaceful event, nothing akin to a supernova, which the Sun is too small to undergo as part of its evolution. Any observer present to witness this occurrence would see a massive increase in the speed of the solar wind, but not enough to destroy a planet completely. However, the star's loss of mass could send the orbits of the surviving planets into chaos, causing some to collide, others to be ejected from the Solar System, and others to be torn apart by tidal interactions. Afterwards, all that will remain of the Sun is a white dwarf, an extraordinarily dense object, 54% of its original mass but only the size of Earth. Initially, this white dwarf may be 100 times as luminous as the Sun is now. It will consist entirely of degenerate carbon and oxygen but will never reach temperatures hot enough to fuse these elements. Thus, the white dwarf Sun will gradually cool, growing dimmer and dimmer.
As the Sun dies, its gravitational pull on orbiting bodies, such as planets, comets, and asteroids, will weaken due to its mass loss. All remaining planets' orbits will expand; if Venus, Earth, and Mars still exist, their orbits will lie roughly at , , and , respectively. They and the other remaining planets will become dark, frigid husks, completely devoid of life. They will continue to orbit their star, their speed slowed due to their increased distance from the Sun and the Sun's reduced gravity. Two billion years later, when the Sun has cooled to the range, the carbon and oxygen in the Sun's core will freeze, with over 90% of its remaining mass assuming a crystalline structure. Eventually, after roughly one quadrillion years, the Sun will finally cease to shine altogether, becoming a black dwarf.
Galactic interaction
The Solar System travels alone through the Milky Way in a circular orbit approximately 30,000 light years from the Galactic Center. Its speed is about 220 km/s. The period required for the Solar System to complete one revolution around the Galactic Center, the galactic year, is in the range of 220–250 million years. Since its formation, the Solar System has completed at least 20 such revolutions.
Various scientists have speculated that the Solar System's path through the galaxy is a factor in the periodicity of mass extinctions observed in the Earth's fossil record. One hypothesis supposes that vertical oscillations made by the Sun as it orbits the Galactic Centre cause it to regularly pass through the galactic plane. When the Sun's orbit takes it outside the galactic disc, the influence of the galactic tide is weaker; as it re-enters the galactic disc, as it does every 20–25 million years, it comes under the influence of the far stronger "disc tides", which, according to mathematical models, increase the flux of Oort cloud comets into the Solar System by a factor of 4, leading to a massive increase in the likelihood of a devastating impact.
However, others argue that the Sun is currently close to the galactic plane, and yet the last great extinction event was 15 million years ago. Therefore, the Sun's vertical position cannot alone explain such periodic extinctions, and that extinctions instead occur when the Sun passes through the galaxy's spiral arms. Spiral arms are home not only to larger numbers of molecular clouds, whose gravity may distort the Oort cloud, but also to higher concentrations of bright blue giants, which live for relatively short periods and then explode violently as supernovae.
Galactic collision and planetary disruption
Although the vast majority of galaxies in the Universe are moving away from the Milky Way, the Andromeda Galaxy, the largest member of the Local Group of galaxies, is heading toward it at about 120 km/s. In 4 billion years, Andromeda and the Milky Way will collide, causing both to deform as tidal forces distort their outer arms into vast tidal tails. If this initial disruption occurs, astronomers calculate a 12% chance that the Solar System will be pulled outward into the Milky Way's tidal tail and a 3% chance that it will become gravitationally bound to Andromeda and thus a part of that galaxy. After a further series of glancing blows, during which the likelihood of the Solar System's ejection rises to 30%, the galaxies' supermassive black holes will merge. Eventually, in roughly 6 billion years, the Milky Way and Andromeda will complete their merger into a giant elliptical galaxy. During the merger, if there is enough gas, the increased gravity will force the gas to the centre of the forming elliptical galaxy. This may lead to a short period of intensive star formation called a starburst. In addition, the infalling gas will feed the newly formed black hole, transforming it into an active galactic nucleus. The force of these interactions will likely push the Solar System into the new galaxy's outer halo, leaving it relatively unscathed by the radiation from these collisions.
It is a common misconception that this collision will disrupt the orbits of the planets in the Solar System. Although it is true that the gravity of passing stars can detach planets into interstellar space, distances between stars are so great that the likelihood of the Milky Way–Andromeda collision causing such disruption to any individual star system is negligible. Although the Solar System as a whole could be affected by these events, the Sun and planets are not expected to be disturbed.
However, over time, the cumulative probability of a chance encounter with a star increases, and disruption of the planets becomes all but inevitable. Assuming that the Big Crunch or Big Rip scenarios for the end of the Universe do not occur, calculations suggest that the gravity of passing stars will have completely stripped the dead Sun of its remaining planets within 1 quadrillion (1015) years. This point marks the end of the Solar System. Although the Sun and planets may survive, the Solar System, in any meaningful sense, will cease to exist.
Chronology
The time frame of the Solar System's formation has been determined using radiometric dating. Scientists estimate that the Solar System is 4.6 billion years old. The oldest known mineral grains on Earth are approximately 4.4 billion years old. Rocks this old are rare, as Earth's surface is constantly being reshaped by erosion, volcanism, and plate tectonics. To estimate the age of the Solar System, scientists use meteorites, which were formed during the early condensation of the solar nebula. Almost all meteorites (see the Canyon Diablo meteorite) are found to have an age of 4.6 billion years, suggesting that the Solar System must be at least this old.
Studies of discs around other stars have also done much to establish a time frame for Solar System formation. Stars between one and three million years old have discs rich in gas, whereas discs around stars more than 10 million years old have little to no gas, suggesting that giant planets within them have ceased forming.
Timeline of Solar System evolution
Note: All dates and times in this chronology are approximate and should be taken as an order of magnitude indicator only.
| Physical sciences | Physical cosmology | null |
6139788 | https://en.wikipedia.org/wiki/Bending%20%28metalworking%29 | Bending (metalworking) | Bending is a manufacturing process that produces a V-shape, U-shape, or channel shape along a straight axis in ductile materials, most commonly sheet metal. Commonly used equipment include box and pan brakes, brake presses, and other specialized machine presses. Typical products that are made like this are boxes such as electrical enclosures and rectangular ductwork.
Process
In press brake forming, the work piece is positioned over a die block and a punch then presses the sheet into the die block to form a shape. Usually bending has to overcome both tensile stresses and compressive stresses. When bending is done, the residual stresses cause the material to towards its original position, so the sheet must be over-bent to achieve the proper bend angle. The amount of spring back is dependent on the material, and the type of forming. When sheet metal is bent, it stretches in length. The bend deduction is the amount the sheet metal will stretch when bent as measured from the outside edges of the bend. The bend radius refers to the inside radius. The formed bend radius is dependent upon the dies used, the material properties, and the material thickness.
The U-punch forms a U-shape with a single punch.
Types
There are three basic types of bending on a press brake, each is defined by the relationship of the end tool position to the thickness of the material. These three are Air Bending, Bottoming and Coining. The configuration of the tools for these three types of bending are nearly identical. A die with a long rail form tool with a radiused tip that locates the inside profile of the bend is called a punch. Punches are usually attached to the ram of the machine by clamps and move to produce the bending force. A die with a long rail form tool that has concave or V-shaped lengthwise channel that locate the outside profile of the form is called a die. Dies are usually stationary and located under the material on the bed of the machine. Note that some locations do not differentiate between the two different kinds of dies (punches and dies). The other types of bending listed use specially designed tools or machines to perform the work.
Air bending
This bending method forms material by pressing a punch (also called the upper or top die) into the material, forcing it into a bottom V-die, which is mounted on the press. The punch forms the bend so that the distance between the punch and the side wall of the V is greater than the material thickness (T).
Either a V-shaped or square opening may be used in the bottom die (dies are frequently referred to as tools or tooling). Because it requires less bend force, air bending tends to use smaller tools than other methods.
Some of the newer bottom tools are adjustable, so, by using a single set of top and bottom tools and varying press-stroke depth, different profiles and products can be produced. Different materials and thicknesses can be bent in varying bend angles, adding the advantage of flexibility to air bending. There are also fewer tool changes, thus, higher productivity.
A disadvantage of air bending is that, because the sheet does not stay in full contact with the dies, it is not as precise as some other methods, and stroke depth must be kept very accurate. Variations in the thickness of the material and wear on the tools can result in defects in parts produced. Thus, the use of adequate process models is important.
Air bending's angle accuracy is approximately ±0.5 deg. Angle accuracy is ensured by applying a value to the width of the V opening, ranging from 6 T (six times material thickness) for sheets to 3 mm thick to 12 T for sheets more than 10 mm thick. Springback depends on material properties, influencing the resulting bend angle.
Depending on material properties, the sheet may be overbent to compensate for springback.
Air bending does not require the bottom tool to have the same radius as the punch. Bend radius is determined by material elasticity rather than tool shape.
The flexibility and relatively low tonnage required by air bending are helping to make it a popular choice. Quality problems associated with this method are countered by angle-measuring systems, clamps and crowning systems adjustable along the x and y axes, and wear-resistant tools.
The K-factor approximations given below are more likely to be accurate for air bending than the other types of bending due to the lower forces involved in the forming process.
Bottoming
In bottoming, the sheet is forced against the V opening in the bottom tool. U-shaped openings cannot be used. Space is left between the sheet and the bottom of the V opening. The optimum width of the V opening is 6 T (T stands for material thickness) for sheets about 3 mm thick, up to about 12 T for 12 mm thick sheets. The bending radius must be at least 0.8 T to 2 T for sheet steel. Larger bend radii require about the same force for bottoming as they do for air bending, however, smaller radii require greater force—up to five times as much—than air bending. Advantages of bottoming include greater accuracy and less springback. A disadvantage is that a different tool set is needed for each bend angle, sheet thickness, and material. In general, air bending is the preferred technique.
Coining
In coining, the top tool forces the material into the bottom die with 5 to 30 times the force of air bending, causing permanent deformation through the sheet. There is little, if any, spring back. Coining can produce an inside radius as low as 0.4 T, with a 5 T width of the V opening. While coining can attain high precision, higher costs mean that it is not often used.
Three-point bending
Three-point bending is a newer process that uses a die with an adjustable-height bottom tool, moved by a servo motor. The height can be set within 0.01 mm. Adjustments between the ram and the upper tool are made using a hydraulic cushion, which accommodates deviations in sheet thickness. Three-point bending can achieve bend angles with 0.25 deg. precision. While three-point bending permits high flexibility and precision, it also entails high costs and there are fewer tools readily available. It is being used mostly in high-value niche markets.
Folding
In folding, clamping beams hold the longer side of the sheet. The beam rises and folds the sheet around a bend profile. The bend beam can move the sheet up or down, permitting the fabricating of parts with positive and negative bend angles. The resulting bend angle is influenced by the folding angle of the beam, tool geometry, and material properties. Large sheets can be handled in this process, making the operation easily automated. There is little risk of surface damage to the sheet.
Wiping
In wiping, the longest end of the sheet is clamped, then the tool moves up and down, bending the sheet around the bend profile. Though faster than folding, wiping has a higher risk of producing scratches or otherwise damaging the sheet, because the tool is moving over the sheet surface. The risk increases if sharp angles are being produced.
This method will typically bottom or coin the material to set the edge to help overcome springback. In this bending method, the radius of the bottom die determines the final bending radius.
Rotary bending
Rotary bending is similar to wiping but the top die is made of a freely rotating cylinder with the final formed shape cut into it and a matching bottom die. On contact with the sheet, the roll contacts on two points and it rotates as the forming process bends the sheet. This bending method is typically considered a "non-marking" forming process suitable to pre-painted or easily marred surfaces. This bending process can produce angles greater than 90° in a single hit on standard press brakes process.
Roll bending
The roll bending process induces a curve into bar or plate workpieces.
There should be proper pre-punching allowance.
Elastomer bending
In this method, the bottom V-die is replaced by a flat pad of urethane or rubber. As the punch forms the part, the urethane deflects and allows the material to form around the punch. This bending method has a number of advantages. The urethane will wrap the material around the punch and the end bend radius will be very close to the actual radius on the punch. It provides a non-marring bend and is suitable for pre-painted or sensitive materials. Using a special punch called a radius ruler with relieved areas on the urethane U-bends greater than 180° can be achieved in one hit, something that is not possible with conventional press tooling. Urethane tooling should be considered a consumable item and while they are not cheap, they are a fraction of the cost of dedicated steel. It also has some drawbacks, this method requires tonnage similar to bottoming and coining and does not do well on flanges that are irregular in shape, that is where the edge of the bent flange is not parallel to the bend and is short enough to engage the urethane pad.
Joggling
Joggling, also known as joggle bending, is an offset bending process in which two opposite bends with equal angles are formed in a single action creating a small s-shape bend profile and an offset between the unbent face and the result flange that is typically less than 5 material thicknesses. Often the offset will be one material thickness, in order to allow a lap joint where the edge of one sheet of material is laid on top of the other.
Calculations
Many variations of these formulas exist and are readily available online. These variations may often seem to be at odds with one another, but they are invariably the same formulas simplified or combined. What is presented here are the unsimplified formulas.
All formulas use the following keys:
Lf = flat length of the sheet
BA = bend allowance
BD = bend deduction
R = inside bend radius
K = K-factor, which is t / T
T = material thickness
t = distance from inside face to the neutral line
A = bend angle in degrees (the angle through which the material is bent)
The neutral line (also called the Neutral axis) is an imaginary profile that can be drawn through a cross-section of the workpiece that represents the locus where no tensile or compressive stress are present but shear stresses are at their maximum. In the bend region, the material between the neutral line and the inside radius will be under compression during the bend while the material between the neutral line and the outside radius will be under tension during the bend. Its location in the material is a function of the forces used to form the part and the material yield and tensile strengths. This theoretical definition also coincides with the geometric definition of the plane representing the unbent flat pattern shape within the cross-section of the bent part. Furthermore, the bend allowance (see below) in air bending depends primarily on the width of the opening of the bottom die. As a result, the bending process is more complicated than it appears to be at first sight.
Both bend deduction and bend allowance represent the difference between the neutral line or unbent flat pattern (the required length of the material prior to bending) and the formed bend. Subtracting them from the combined length of both flanges gives the flat pattern length. The question of which to use is determined by the dimensioning method used to define the flanges as shown in the two diagrams below. The flat pattern length is always shorter in length than the sum of all the flange length dimensions due to the geometric transformation. This gives rise to the common perspective that that material is stretching during bending and the bend deduction and bend allowance are the distance that each bend stretches. While a helpful way to look at it, a careful examination of the formulas and stresses involved show this to be false.
Most 3D Solid Modeling CAD software has sheet metal functions or add-ons that performs these calculations automatically.
Bend allowance
The bend allowance (BA) is the length of the arc of the neutral line between the tangent points of a bend in any material. Adding the length of each flange as dimensioned by B in the diagram to the BA gives the Flat Pattern length. This bend allowance formula is used to determine the flat pattern length when a bend is dimensioned from 1) the center of the radius, 2) a tangent point of the radius (B) or 3) the outside tangent point of the radius on an acute angle bend (C). When dimensioned to the outside tangent, the material thickness and bend radius are subtracted from it to find the dimension to the tangent point of the radius before adding in the bend allowance.
The BA can be estimated using the following formula, which incorporates the empirical K-factor:
Bend deduction
The bend deduction BD is defined as the difference between the sum of the flange lengths (from the edge to the apex) and the initial flat length.
The outside set back (OSSB) is the length from the tangent point of the radius to the apex of the outside of the bend. The bend deduction (BD) is twice the outside setback minus the bend allowance. BD is calculated using the following formula, where A is the angle in radians (=degrees*π/180):
For bends at 90 degrees this formula can be simplified to:
K-factor
K-factor is a ratio of the location of the neutral line to the material thickness as defined by t/T where t = location of the neutral line and T = material thickness. The K-factor formula does not take the forming stresses into account but is simply a geometric calculation of the location of the neutral line after the forces are applied and is thus the roll-up of all the unknown (error) factors for a given setup. The K-factor depends on many variables including the material, the type of bending operation (coining, bottoming, air-bending, etc.) the tools, etc. and is typically between 0.3 and 0.5.
The following equation relates the K-factor to the bend allowance:
The following table is a "rule of thumb". Actual results may vary remarkably.
The following formula can be used in place of the table as a good approximation of the K-factor for air bending:
Advantages and disadvantages
Bending is a cost-effective near net shape process when used for low to medium quantities. Parts usually are lightweight with good mechanical properties. A disadvantage is that some process variants are sensitive to variations in material properties. For instance, differences in spring-back have a direct influence on the resulting bend angle. To mitigate this, various methods for in-process control have been developed. Other approaches include combining brakeforming with incremental forming.
Broadly speaking, each bend corresponds with a set-up (although sometimes, multiple bends can be formed simultaneously). The relatively large number of set-ups and the geometrical changes during bending make it difficult to address tolerances and bending errors a priori during set-up planning, although some attempts have been made
| Technology | Metallurgy | null |
22161339 | https://en.wikipedia.org/wiki/Myxogastria | Myxogastria | Myxogastria/Myxogastrea (myxogastrids, ICZN) or Myxomycetes (ICN) is a class of slime molds that contains 5 orders, 14 families, 62 genera, and 888 species. They are colloquially known as the plasmodial or acellular slime moulds.
All species pass through several very different morphologic phases, such as microscopic individual cells, slimy amorphous organisms visible with the naked eye, and conspicuously shaped fruit bodies. Although they are monocellular, they can reach immense widths and weights: in extreme cases they can be up to across and weigh up to .
The class Myxogastria is distributed worldwide, but it is more common in temperate regions where it has a higher biodiversity than in polar regions, the subtropics, or the tropics. They are mainly found in open forests, but also in extreme regions such as deserts, under snow blankets, or underwater. They also occur on the bark of trees, sometimes high in the canopy. These are known as corticolous myxomycetes. Most species are very small.
Taxonomy and classification
Nomenclature
Myxomycota, now considered a synonym of Myxogastria, comes from the Ancient Greek words μύξα , which means "mucus", and μύκης , which means "fungus". The name Myxogastria was introduced in 1970 by Lindsay Shepherd Olive to describe the family Myxogastridae, which was introduced in 1899 by Thomas Huston Macbride. Swedish mycologist Elias Magnus Fries described numerous slime moulds as Myxogasteres in 1829. Species in the class Myxogastria are colloquially known as plasmodial or acellular slime moulds. Some consider the Myxogastria a separate kingdom, with an unsettled phylogeny because of conflicting molecular and developmental data. The relations among myxogastrid orders are as yet unclear.
Range
The continuous classification of new taxa reveals that the class is not fully described. According to a 2000 inquiry, there were 1012 officially accepted taxa, including 866 on species level. Another study in 2007 stated a number of more than 1000, in which the Myxogastria comprised the biggest group of slime moulds, with over 900 species. On the basis of sequenced environmental samples it is estimated that the group has between 1200 and 1500 species – more than previously estimated. Among the 1012 taxa only a few species are common: 305 species were discovered in a single location or groupings, a further 258 species were found in a few areas between 2–20 times, and only 446 were common in several locations with over 20 discoveries.
Reclassifications encounter problems because the Myxogastriae are morphologically very plastic, which is to say susceptible to environmental influences; only a few characteristics are diagnostic for a small number of species. In the past, authors have unsuccessfully tried to describe a new taxon based on a small number of examples, but this leads to numerous duplications, sometimes even at genus level. For example, Squamuloderma nullifila is actually a species from the genus Didymium.
Classification and phylogeny
The following classification is based on Adl et al. (2005) while the classes and further divisions on Dykstra & Keller (2000), who included the Myxogastria in Mycetozoa. The sister taxon is the subclass Dictyostelia. Together with the Protostelia they formed the taxon Eumycetozoa. Other subclasses differ from the other species mainly in the development of fruit bodies; while Protostelia create a separate fruit body from each single mononuclear cell, Dictyostelia develop cell complexes – the so-called pseudo-plasmodia – from separate cells, which then become fruit bodies.
Clade Myxogastria (or myxomycetes)
Class Ceratiomyxomycetes Hawksworth, Sutton & Ainsworth 1983
Order Ceratiomyxida Martin 1949 ex Farr & Alexopoulos 1977
Order Protosporangiida Shadwick & Spiegel 2012
Class Myxomycetes Link 1833 em. Haeckel 1866
Subclass Lucisporomycetidae Leontyev et al. 2019 (Clear‑spored acellular slime moulds)
Superorder Cribrarianae Leontyev 2015
Order Cribrariales Macbride 1922
Family Cribrariaceae Rostafinski 1873
Superorder Trichianae Leontyev 2015
Order Reticulariales Leontyev 2015
Family Reticulariaceae Rostafinski 1873
Order Liceales Jahn 1928
Family Liceaceae Rostafinski 1873
Order Trichiales Macbride 1922
Family Dianemidae Macbride 1899
Family Trichiidae Rostafinski 1826
Subclass Collumellidia Leontyev et al. 2019 (Dark‑spored acellular slime moulds)
Order Echinosteliopsidales Shchepin et al.
Family Echinosteliopsidaceae Olive 1970
Superorder Echinostelianae Leontyev 2015
Order Echinosteliales Martin 1961
Family Echinosteliaceae Rostafinski ex Cooke 1877
Superorder Stemonitanae Leontyev 2015
Order Clastodermatales Leontyev et al. 2019
Family Clastodermataceae Alexopoulos & Brooks 1971
Order Meridermatales Leontyev 2015
Family Meridermataceae Leontyev 2015
Order Stemonitales Macbride 1922
Family Comatrichaceae Leontyev 2015
Family Stemonitidaceae Fries 1829
Order Physarales Macbride 1922
Family Didymiaceae Rostafinski ex Cooke 1877
Family Lamprodermataceae Leontyev 2015
Family Physaraceae Chevallier 1826
Some classifications place part of the orders above in the subclass Myxogastromycetidae.
Characteristics and life cycle
Monocellular, mononuclear phase
Spores
The spores of Myxogastria are haploid, mainly round and measure between 5 μm and 20 μm, rarely up to 24 μm in diameter. Their surface is generally reticular, sharp, warty or spiky and very rarely smooth. The typical colour of the spore mass becomes visible through the structure, since the spores themselves are not pigmented. In some species, especially of the genus Badhamia, the spores produce lumps. The colour, shape and diameter of spores are important characteristics for identifying species.
Important factors for the germination of spores are mainly moisture and temperature. The spores usually remain germinable after several years; there were even spores preserved in herbarium specimens which germinated after 75 years. After the spores' development, they first receive a diploid nucleus, and the meiosis takes place in the spore. At the germination, the spore shells open either alongside special germinal pores or chinks, or rip irregularly and then release one to four haploid protoplasts.
Myxamoebae and Myxoflagellates
In those species which reproduce sexually, haploid cells bud from the spores. Depending on the environmental conditions, either a myxamoeba or a myxoflagellate buds from the spore. Myxamoebae move like amoebae – that is, crawling on the substrate – and are produced in dry conditions. Myxoflagellates, which are peritrichous and can swim, develop in moist to wet environments. Myxoflagellates almost always have two flagella; one is generally shorter than the other and sometimes only vestigial. The flagella are used for locomotion and to help to move food particles closer. If the humidity changes, cells can switch between the two manifestations. Neither form has a cell wall. This developmental stage (and the next one) serves as a nourishment provider and is also known as the first trophic phase (nourishment phase). In this monocellular phase, the Myxogastria consume bacteria and fungus spores, and probably dissolved substances, and they reproduce through simple cell division. If the environmental conditions change adversely in this phase, for example extreme temperature, extreme dryness or food shortage, the Myxogastria may switch to very long-lived, thin-shelled quiescent states – the so-called microcysts. For that to happen, the myxamoebae assume a round shape and secrete a thin cell wall. In this state they can easily survive one year or longer. If living conditions improve, they become active again.
Zygogenesis
If two cells of the same type meet in this phase, they cross-fertilise to a diploid zygote through the fusion of protoplasms and nuclei. The conditions which trigger this are not known. The diploid zygote becomes a multinucleated plasmodium through multiple nuclear divisions without further cell division. If the resulting cells were peritrichous, they change their shape before the fusion from the peritrichous form to the myxamoeba. The production of a zygote requires two cells of different mating types (heterothallic).
Plasmodium
The second trophic phase begins with the development of the plasmodium. The multinucleated organism now absorbs via phagocytosis as many nutrients as possible. These are bacteria, protists, dissolved substances, moulds, higher fungi and small particles of organic material. This enables the cell to undergo enormous growth. The nucleus divides multiple times, and the cell soon becomes visible to the naked eye and usually has a surface area – depending on the species – up to one square metre; however, in 1987 one artificially cultivated cell of Physarum polycephalum attained a surface area of 5.5 sq m. Myxogastria species have numerous nuclei in their trophic plasmodium phase; the small, non-veined proto-plasmodia have between 8–100 nuclei, while large, veined meshworks have between 100 and 10 million nuclei. All of these remain part of a single cell, which has a viscous, slimy consistency, and may be transparent, white, or brightly coloured in orange, yellow, or pink.
The cell has chemotactic and negative phototactic capabilities in this phase, meaning that it is able to move towards nutrients and away from dangerous substances and light. The movements originate in the grainy cytoplasm, which streams by pulsation in one direction within the cell. In this way the cell reaches a speed of up to 1000 μm per second – the speed in plant cells is 2–78 μm per second. A resting state, the so-called sclerotium, may occur in this phase. The sclerotium is a hardened, resistant form composed of numerous "macrocysts", which enable the myxogastria to survive in adverse conditions, for example during winter or dry periods, in this phase.
Fructification
Mature plasmodia can produce fruit bodies under appropriate circumstances. The exact triggers for this process are unknown. According to laboratory researchers, changes in humidity, temperature or pH value as well as starvation periods were thought to be the triggers in some species. The plasmodia abandon their nutrient intake and crawl, attracted by light – a positive phototaxis – towards a dry, light area, to get an optimal spread of the spores. Once the fructification begins, it cannot be stopped. If disturbances occur, malformed spore-bearing fruit bodies are often produced.
The plasmodium or parts of the fruit bodies can be smaller than one millimetre, in extreme cases they are up to a square metre and weigh up to (Brefeldia maxima). Their shape is often pediculated or unstiped sporangia with non-cellular stems, but can also appear as veined or netted plasmodiocarps, pincushion-shaped aethaliae or seemingly pincushion-shaped pseudo-aethaliae. The fruit bodies almost always have a hypothallus on the edge. The abundantly produced spores are stored in a reticular or filamentous structure – the so-called capillitium – and are found on nearly all species except Liceida and other species from the genus Echinostelium. When the open fruit bodies have dried, the spores are dispersed by wind or by small animals such as woodlice, mites or beetles, which either pick up the spores through contact with the fruit bodies or ingest and then excrete them. Dispersal by running water is also possible, but it plays a minor role.
Asexual forms
Some Myxogastria species may produce asexually. These are continuously diploid. There is no meiosis before the germination of the spores and the production of the plasmodium proceeds without germination of two cells.
Distribution and ecology
Distribution
Myxogastria are distributed worldwide; species were found by early researchers on all continents. However, as many parts of the world were yet not discovered or explored, the exact distribution is not fully known. Europe and North America are often considered the basic habitat of the Myxogastria species. According to recent research, the majority of species are not widely distributed. The Myxogastria are most commonly found in temperate latitudes, and rarely in the polar regions, the subtropics or tropics. The physical features of the substrate and climatic conditions are the major aspects of the species' presence. Endemism is rare.
In the northern areas, the species can be found in Alaska, Iceland, northern Scandinavia, Greenland and Russia. These are not only particular, specialised species; according to an overview study, more than 150 species were found in the arctic and subarctic regions of Iceland, Greenland, northern Russia and Alaska. These distinctly exceed the tree line. In Greenland, the habitat may reach the 77th latitude line. The Myxogastria species reach their largest biodiversity and highest frequency in forests of temperate regions, which are ideal habitats because of the amount of rich organic material, suitable humidity (not too high) and long-lasting snow cover for snow-inhabiting species.
Few Myxogastria species are found in the tropics and subtropics, mainly because of the high humidity which prevents the necessary dehydration of the fruit bodies to permit spore dispersal and promotes infestation by moulds. Other factors are low light levels under the forest canopy which reduces phototaxis, light winds, poor soils, natural enemies and heavy rainfall which can wash away or destroy cells. Species living in soil or deadwood decrease as humidity increases. In a study from Costa Rica, 73% of the total findings were in the relatively dry Tropical Moist Forest, while 18% were in the very moist Tropical Premontane Wet Forests and only 9% in Lower Premontane Rain Forest.
In the Antarctic, species were found in the South Shetland Islands, South Orkney Islands, South Georgia and the Antarctic Peninsula. Species from the Antarctic or subantarctic regions are rarer than specimens in the Arctic regions, although lack of access may be a factor. Until 1983, only 5 records were made, with only individual finds since then. According to two studies of the myxomycete flora of these regions, more species were discovered in the subantarctic forests, for example 67 species in Argentinian Patagonia and Tierra del Fuego, and 22 on high ground on Macquarie Island.
Habitats
The majority of Myxogastria species live terrestrially in open forests. The most important microhabitat is deadwood, but also the bark of living trees (corticolous myxomycetes), rotting plant material, soil, and animal excrements. Slime moulds may be found in numerous unusual locations. The comprehensive group of the nivicol Myxogastria populate closed snow blankets, to quickly fructificate at exposure – for example during thaws – and release their spores. Other habitats are deserts – 33 species were found in the Sonora desert, for example – or living on leaves from plants in the tropics. Some species live in aquatic environments, such as those of the genera Didymium, Physarum, Perichaena, Fuligo, Comatricha and Licea, which were found living underwater as myxoflagelletes and plasmodia. All but one species, Didymium difforme, fructificated only when the water ebbed or when they left it.
Relationship to other creatures
The relationships of the Myxogastria to other creatures have not been thoroughly researched as of 2012. Their natural predators include many arthropods, including mites and springtails, and especially beetles such as the rove beetles, round fungus beetles, wrinkled bark beetles, Eucinetidae, Clambidae, Eucnemidae (false click beetles), Sphindidae, Cerylonidae, and minute brown scavenger beetles. Various Nematodes have also been observed to be their predators; they attach their posterior portion on the cytosol of the plasmodia or even live within the strands. Certain Diptera species have evolved to specialise in this way: these are mostly representatives of the Mycetophilidae, Sciaridae and Drosophilidae. The species Epicypta testata was especially frequently found, especially on Enteridium lycoperdon, Enteridium splendens, Lycogala epidendrum, and Tubifera ferruginosa.
Some true fungi specialise in the colonisation of the Myxogastriae: almost all of these are species of sac fungi. The most common such fungus is Verticillium rexianum – mainly species from Comatricha or Stemonitis. Gliocladium album and Sesquicillium microsporum are often found on Physaridae, while Polycephalomyces tomentosus is often found on certain species of Trichiidae. Nectriopsis violacea specialises on Fuligo septica. Bacterial associates, mainly from the family Enterobacteriaceae, were discovered on plasmodia. The combination of plasmodia and bacteria can bind atmospheric nitrogen or produce enzymes which make possible the decomposition of e.g. lignin, carboxymethylcellulose, or xylan. In a few cases, the plasmodia acquired salt tolerance or tolerance of heavy metals through this association.
Some myxomycetes (Physarum) cause disease in plants such as turfgrasses, but no control is usually necessary against them.
Fossil records
Fossil records of Myxogastria are extremely rare. Due to their short lifespan and the fragile structures of the plasmodia and the fruit body, fossilisation and similar processes are not possible. Only their spores can be mineralised. The few known examples of fossilised living states are preserved in amber.
three fruit bodies, two spores and one plasmodium have been described. Two older taxa – Charles Eugène Bertrand's Myxomycetes mangini and Bretonia hardingheni from 1892 – are now considered dubious and are today often disregarded. Friedrich Walter Domke described in 1952 a 35 to 40 million year old find in Baltic amber of Stemonitis splendens, an extant species. The state and completeness of the fruit bodies are remarkable, enabling accurate determination. From the same period, location and material is an Arcyria sulcata, first described in 2003 by Heinrich Dörfelt and Alexander Schmidt, a species very similar to today's Arcyria denudata. Both discoveries imply that the fruit bodies of the Myxogastria have changed only slightly in the last 35–40 million years.
However, the Protophysarum balticum from Baltic amber, first described by Dörfelt and Schmidt in 2006, is considered questionable. The fossil was inconsistent with the typical characteristics of the genus and it was not a valid publication because no Latin name was identified with it. Also, important details of its fruit bodies were not visible or contradicted the identification. Today it is assumed that the fossil belongs to a lichen similar to the genus Chaenotheca.
The only known discovery of a preserved plasmodium was found in Dominican amber, and was then grouped into the Physarida. However, this claim is also considered doubtful as the publication was later classified as insufficient due to lack of evidence.
In 2019 sporocarps belonging to Stemonitis was described from Burmese amber, considered to be of a mid-Cretaceous age around 99 million years old. The sporocarps are indistinguishable from extant taxa, suggesting a long morphological stasis.
The only known mineralised fossils are the two spore findings from 1971, one of which, Trichia favoginea, is assumed to be from the postglacial period. In palynologian researches, by absorbing Myxogastria spores, the fossil was not recognised.
History of research
Because of their unprepossessing nature, the Myxogastriae were for a long time not well researched. Thomas Panckow first named the mould Lycogala epidendrum as "Fungus cito crescentes" (fast-growing fungus) in his 1654 book Herbarium Portatile, oder behendes Kräuter- und Gewächsbuch. In 1729, Pier Antonio Micheli thought that fungi are different from moulds, and Heinrich Friedrich Link agreed with this hypothesis in 1833. Elias Magnus Fries documented the plasmodial stage in 1829, and 35 years later Anton de Bary observed the germination of the spores. De Bary also discovered the cyclosis in the cell for the movement, he saw them as animal-like creatures and reclassified them as Mycetozoa, which literally translates "Fungus animals". This interpretation prevailed until the second half of the 20th century.
From 1874 to 1876, Józef Tomasz Rostafiński, a student of Anton de Bary, published the first extensive monograph on the group. Three monographs by Arthur Lister and Guilielma Lister were published in 1894, 1911, and 1925. These were groundbreaking works about the Myxogastria, as was the 1934 book The Myxomycetes by Thomas H. Macbride and George Willard Martin. Important works in the late 20th century were the 1969 monographs by George Willard Martin and Constantine John Alexopoulos, and the 1975 monograph by Lindsay Shepherd Olive. The first is perhaps the most notable, as with it "the modern era of the taxonomy of the Myxogastria began". Other notable researchers were Persoon, Rostafinski, Lister, Macbridge, and Martin and Alexopoulos, who discovered and classified many species.
| Biology and health sciences | Eukaryotes | Plants |
10432749 | https://en.wikipedia.org/wiki/Pacific%20bluefin%20tuna | Pacific bluefin tuna | The Pacific bluefin tuna (Thunnus orientalis) is a predatory species of tuna found widely in the northern Pacific Ocean, but it is migratory and also recorded as a visitor to the south Pacific.
In the past it was often included in T. thynnus, the 'combined' species then known as the northern bluefin tuna (when treated as separate, T. thynnus is called the Atlantic bluefin tuna). It may reach as much as in length and in weight.
Like the closely related Atlantic bluefin and southern bluefin, the Pacific bluefin is a commercially valuable species and several thousand tonnes are caught each year. It was considered overfished and subject to overfishing for decades, but catches were reduced in 2011 in order to rebuild the stock and a 2024 stock assessment determined that the species had rebuilt and was no longer overfished nor subject to overfishing. It is now considered a management success. Monterey Bay Aquarium's Seafood Watch program lists Pacific bluefin tuna as a "Good alternative".
Distribution
The Pacific bluefin tuna is primarily found in the North Pacific, ranging from the East Asian coast to the western coast of North America. It is mainly a pelagic species found in temperate oceans, but it also ranges into the tropics and more coastal regions. It typically occurs from the surface to , but has been recorded as deep as .
It spawns in the northwestern Philippine Sea (e.g., off Honshu, Okinawa and Taiwan) and in the Sea of Japan. Some of these migrate to the East Pacific and return to the spawning grounds after a few years. It has been recorded more locally as a visitor to the Southern Hemisphere, including off Australia, New Zealand, the Gulf of Papua and French Polynesia.
The species is considered to consist of only one stock.
Physiology
Thermoregulation
Almost all fish are cold-blooded (ectothermic). However, tuna and mackerel sharks are warm-blooded: they can regulate their body temperature. Warm-blooded fish possess organs near their muscles called retia mirabilia that consist of a series of minute parallel veins and arteries that supply and drain the muscles. As the warmer blood in the veins returns to the gills for fresh oxygen it comes into close contact with cold, newly oxygenated blood in the arteries. The system acts as a counter-current heat exchanger and the heat from the blood in the veins is given up to the colder arterial blood rather than being lost at the gills. The net effect is less heat loss through the gills. Fish from warmer water elevate their temperature a few degrees whereas those from cold water may raise it as much as warmer than the surrounding sea.
The tuna's ability to maintain body temperature has several definite advantages over other sea life. It need not limit its range according to water temperature, nor is it dominated by climatic changes. The additional heat supplied to the muscles is also advantageous because of the resulting extra power and speed.
Life cycle
Pacific bluefin tunas reach maturity at about 5 years of age, the generation length is estimated at 7–9 years and based on two separate sources the longevity is 15 years or 26 years. At maturity it is about long and weighs about . Individuals that are long are regularly seen, and the maximum reported is in length and in weight. Elsewhere, a mass of up to has been reported for the species. According to the International Game Fish Association, the all-tackle game fish record was a individual (Donna Pascoe) caught on 19 February 2014 onboard charter boat Gladiator during the National Tournament.
Spawning occurs from April to August, but the exact timing depends on the region: Early in the northwest Philippine Sea (the southern part of its breeding range) and late in the Sea of Japan (the northern part of its breeding range). Large females can carry more eggs than small ones, and between 5 million and 25 million eggs have been reported.
Pacific bluefins eat various small schooling squids and fishes, but have also been recorded taking sessile animals, pelagic red crabs and krill.
Human interaction
Commercial fishery
Pacific bluefin tuna support a large commercial fishery.
Aquaculture
Japan is both the biggest consumer and the leader in tuna farming research. Kinki University of Japan first successfully farmed already-hatched bluefin tuna in 1979. In 2002, they succeeded in breeding them, and in 2007, the process was repeated for a third generation. This farm-raised tuna is now known as Kindai tuna. Kindai is a contraction of Kinki University (Kinki daigaku).
Conservation
Unlike the other bluefins (Atlantic and southern), the Pacific bluefin tuna was not considered threatened initially, resulting in a Least Concern rating in 2011. In 2014, it was found to be threatened and the status was changed to Vulnerable. The current status is listed as Near threatened. Based on a 2024 stock assessment, it was considered to have been rebuilt and not overfished, nor subject to overfishing.
According to the 2024 stock assessment by the International Scientific Committee for Tuna and Tuna-Like Species in the North Pacific Ocean (ISC), the population has increased from a low point of about 2 percent of historic levels in 2010 to about 23 percent in 2020. This has coincided with a reduction in fishing mortality due to stricter management measures. The IUCN classifies the population as "Near Threatened", although that designation has not been updated since the stock was found to have been rebuilt.
Catches have ranged between about 8,000 and 40,000 tonnes since 1952.
Its wide range and migratory behavior lead to some problems, since fisheries in the species are managed by several different Regional Fisheries Management Organisations that have sometimes given conflicting advice. The IUCN have recommended that the responsibility be moved to a single organisation. Other recommendations include a substantial reduction of fishing of this species, especially juveniles. In the past, as much as 90% of the caught Pacific bluefins are juveniles.
Monterey Bay Aquarium's Seafood Watch program lists Pacific bluefin tuna as a "Good alternative".
Mercury levels
Pacific bluefin flesh may contain levels of mercury or PCBs that are harmful to humans who consume it. A similar problem exists in other tuna species.
Cuisine
About 80% of the Pacific and Atlantic bluefin tunas are consumed in Japan, and tunas that are particularly suited for sashimi and sushi can fetch very high prices. The fatty belly meat is known as toro, and prized by sushi chefs. In Japan, some foods made available for the first time of the year are considered good luck, especially bluefin tuna. Winning these new year auctions is often used as a way to get publicity, which raises the prices considerably higher than their usual market value: on 5 January 2013, a Pacific bluefin tuna caught off northeastern Japan was sold in the first auction of the year at the Tsukiji fish market in Tokyo for a record 155.4 million yen (US$1.76 million) – leading to record unit prices of US$3,603 per pound, or ¥703,167 per kilogram. A pacific bluefin tuna sold for 333.6 million yen (US$3.1 million) at a Tokyo's Toyosu fish market on 5 January 2019. The price equates to roughly $5,000 a pound, close to double the previous record. The fish was caught off Oma in northern Japan.
| Biology and health sciences | Acanthomorpha | Animals |
12743856 | https://en.wikipedia.org/wiki/Polymer%20engineering | Polymer engineering | Polymer engineering is generally an engineering field that designs, analyses, and modifies polymer materials. Polymer engineering covers aspects of the petrochemical industry, polymerization, structure and characterization of polymers, properties of polymers, compounding and processing of polymers and description of major polymers, structure property relations and applications.
History
The word “polymer” was introduced by the Swedish chemist J. J. Berzelius. He considered, for example, benzene (C6H6) to be a polymer of ethyne (C2H2). Later, this definition underwent a subtle modification.
The history of human use of polymers has been long since the mid-19th century, when it entered the chemical modification of natural polymers. In 1839, Charles Goodyear found a critical advance in the research of rubber vulcanization, which has turned natural rubber into a practical engineering material. In 1870, J. W. Hyatt uses camphor to plasticize nitrocellulose to make nitrocellulose plastics industrial. 1907 L. Baekeland reported the synthesis of the first thermosetting phenolic resin, which was industrialized in the 1920s, the first synthetic plastic product. In 1920, H. Standinger proposed that polymers are long-chain molecules that are connected by structural units through common covalent bonds. This conclusion laid the foundation for the establishment of modern polymer science. Subsequently, Carothers divided the synthetic polymers into two broad categories, namely a polycondensate obtained by a polycondensation reaction and an addition polymer obtained by a polyaddition reaction. In the 1950s, K. Ziegler and G. Natta discovered a coordination polymerization catalyst and pioneered the era of synthesis of stereoregular polymers. In the decades after the establishment of the concept of macromolecules, the synthesis of high polymers has achieved rapid development, and many important polymers have been industrialized one after another.
Classification
The basic division of polymers into thermoplastics, elastomers and thermosets helps define their areas of application.
Thermoplastics
Thermoplastic refers to a plastic that has heat softening and cooling hardening properties. Most of the plastics we use in our daily lives fall into this category. It becomes soft and even flows when heated, and the cooling becomes hard. This process is reversible and can be repeated. Thermoplastics have relatively low tensile moduli, but also have lower densities and properties such as transparency which make them ideal for consumer products and medical products. They include polyethylene, polypropylene, nylon, acetal resin, polycarbonate and PET, all of which are widely used materials.
Elastomers
An elastomer generally refers to a material that can be restored to its original state after removal of an external force, whereas a material having elasticity is not necessarily an elastomer. The elastomer is only deformed under weak stress, and the stress can be quickly restored to a polymer material close to the original state and size. Elastomers are polymers which have very low moduli and show reversible extension when strained, a valuable property for vibration absorption and damping. They may either be thermoplastic (in which case they are known as Thermoplastic elastomers) or crosslinked, as in most conventional rubber products such as tyres. Typical rubbers used conventionally include natural rubber, nitrile rubber, polychloroprene, polybutadiene, styrene-butadiene and fluorinated rubbers.
Thermosets
A thermosetting resin is used as a main component, and a plastic which forms a product is formed by a cross-linking curing process in combination with various necessary additives. It is liquid in the early stage of the manufacturing or molding process, and it is insoluble and infusible after curing, and it cannot be melted or softened again. Common thermosetting plastics are phenolic plastics, epoxy plastics, aminoplasts, unsaturated polyesters, alkyd plastics, and the like. Thermoset plastics and thermoplastics together constitute the two major components of synthetic plastics. Thermosetting plastics are divided into two types: formaldehyde cross-linking type and other cross-linking type.
Thermosets includes phenolic resins, polyesters and epoxy resins, all of which are used widely in composite materials when reinforced with stiff fibers such as fiberglass and aramids. Since crosslinking stabilises the thermoset polymer matrix of these materials, they have physical properties more similar to traditional engineering materials like steel. However, their very much lower densities compared with metals makes them ideal for lightweight structures. In addition, they suffer less from fatigue, so are ideal for safety-critical parts which are stressed regularly in service.
Materials
Plastic
Plastic is a polymer compound which is polymerized by polyaddition polymerization and polycondensation. It is free to change the composition and shape. It is made up of synthetic resins and fillers, plasticizers, stabilizers, lubricants, colorants and other additives. The main component of plastic is resin. Resin means that the polymer compound has not been added with various additives. The term resin was originally named for the secretion of oil from plants and animals, such as rosin and shellac. Resin accounts for approximately 40% - 100% of the total weight of the plastic. The basic properties of plastics are mainly determined by the nature of the resin, but additives also play an important role. Some plastics are basically made of synthetic resins, with or without additives such as plexiglass, polystyrene, etc.
Fiber
Fiber refers to a continuous or discontinuous filament of one substance. Animals and plant fibers play an important role in maintaining tissue. Fibers are widely used and can be woven into good threads, thread ends and hemp ropes. They can also be woven into fibrous layers when making paper or feel. They are also commonly used to make other materials together with other materials to form composites. Therefore, whether it is natural or synthetic fiber filamentous material. In modern life, the application of fiber is ubiquitous, and there are many high-tech products.
Rubber
Rubber refers to highly elastic polymer materials and reversible shapes. It is elastic at room temperature and can be deformed with a small external force. After removing the external force, it can return to the original state. Rubber is a completely amorphous polymer with a low glass transition temperature and a large molecular weight, often greater than several hundred thousand. Highly elastic polymer compounds can be classified into natural rubber and synthetic rubber. Natural rubber processing extracts gum rubber and grass rubber from plants; synthetic rubber is polymerized by various monomers. Rubber can be used as elastic, insulating, water-impermeable air-resistant materials.
Applications
Polyethylene
Commonly used polyethylenes can be classified into low density polyethylene (LDPE), high density polyethylene (HDPE), and linear low density polyethylene (LLDPE). Among them, HDPE has better thermal, electrical and mechanical properties, while LDPE and LLDPE have better flexibility, impact properties and film forming properties. LDPE and LLDPE are mainly used for plastic bags, plastic wraps, bottles, pipes and containers; HDPE is widely used in various fields such as film, pipelines and daily necessities because its resistance to many different solvents.
Polypropylene
Polypropylene is widely used in various applications due to its good chemical resistance and weldability. It has lowest density among commodity plastics. It is commonly used in packaging applications, consumer goods, automatic applications and medical applications. Polypropylene sheets are widely used in industrial sector to produce acid and chemical tanks, sheets, pipes, Returnable Transport Packaging (RTP), etc. because of its properties like high tensile strength, resistance to high temperatures and corrosion resistance.
Composites
Typical uses of composites are monocoque structures for aerospace and automobiles, as well as more mundane products like fishing rods and bicycles. The stealth bomber was the first all-composite aircraft, but many passenger aircraft like the Airbus and the Boeing 787 use an increasing proportion of composites in their fuselages, such as hydrophobic melamine foam. The quite different physical properties of composites gives designers much greater freedom in shaping parts, which is why composite products often look different from conventional products. On the other hand, some products such as drive shafts, helicopter rotor blades, and propellers look identical to metal precursors owing to the basic functional needs of such components.
Biomedical applications
Biodegradable polymers are widely used materials for many biomedical and pharmaceutical applications. These polymers are considered very promising for controlled drug delivery devices. Biodegradable polymers also offer great potential for wound management, orthopaedic devices, dental applications and tissue engineering. Not like non biodegradable polymers, they won't require a second step of a removal from body. Biodegradable polymers will break down and are absorbed by the body after they served their purpose. Since 1960, polymers prepared from glycolic acid and lactic acid have found a multitude of uses in the medical industry. Polylactates (PLAs) are popular for drug delivery system due to their fast and adjustable degradation rate.
Membrane technologies
Membrane techniques are successfully used in the separation in the liquid and gas systems for years, and the polymeric membranes are used most commonly because they have lower cost to produce and are easy to modify their surface, which make them suitable in different separation processes. Polymers helps in many fields including the application for separation of biological active compounds, proton exchange membranes for fuel cells and membrane contractors for carbon dioxide capture process.
Related Major
Petroleum / Chemical / Mineral / Geology
Raw materials and processing
New energy
Automobiles and spare parts
Other industries
Electronic Technology / Semiconductor / Integrated Circuit
Machinery / Equipment / Heavy Industry
Medical equipment / instruments
| Technology | Disciplines | null |
20686278 | https://en.wikipedia.org/wiki/Lobaria%20pulmonaria | Lobaria pulmonaria | Lobaria pulmonaria is a large epiphytic lichen consisting of an ascomycete fungus and a green algal partner living together in a symbiotic relationship with a cyanobacterium—a symbiosis involving members of three kingdoms of organisms. Commonly known by various names like tree lungwort, lung lichen, lung moss, lungwort lichen, oak lungs or oak lungwort, it is sensitive to air pollution and is also harmed by habitat loss and changes in forestry practices. Its population has declined across Europe and L. pulmonaria is considered endangered in many lowland areas. The species has a history of use in herbal medicines, and recent research has corroborated some medicinal properties of lichen extracts.
Description
It is a foliose lichen and its leaf-like thallus is green, leathery and lobed with a pattern of ridges and depressions on the upper surface. Bright green under moist conditions, it becomes brownish and papery when dry. This species often has a fine layers of hairs, a tomentum, on its lower surface. The cortex, the outer protective layer on the thallus surface, is roughly comparable to the epidermis of a green plant. The thallus is typically in diameter, with individual lobes wide and up to 7 cm long. The asexual reproductive structures soredia and isidia are present on the thallus surface. Minute (0.5–1.5 mm in diameter) cephalodia—pockets of cyanobacteria—are often present on the lower surface of the thallus; these spots are conspicuously darker than the green surface of the thallus. Like other foliose lichens, the thallus is only loosely attached to the surface on which it grows.
Photobionts
The thallus contains internal structures known as cephalodia, characteristic of three-membered lichen symbioses involving two photobionts (the photosynthetic symbionts in the lichen relationship). These internal cephalodia, found between the "ribs" of the thallus surface, arise when blue-green algae (from the genus Nostoc) on the thallus surface are enveloped during mycobiont growth. Structurally, cephalodia consist of dense aggregates of Nostoc cells surrounded by thin-walled hyphae—this delimits them from the rest of the thallus which contains a loose structure of thick-walled hyphae. Blue-green cyanobacteria can fix atmospheric nitrogen, providing a nutrient for the lichen. The other photobiont of L. pulmonaria is the green alga Dictyochloropsis reticulata.
Reproduction
Lobaria pulmonaria has the ability to form both vegetative propagation and sexual propagules at an age of about 25 years. In sexual reproduction, the species produces small reddish-brown discs known as apothecia containing asci, from which spores are forcibly released into the air (like ballistospores). Based on studies of ascospore germination, it has been suggested that L. pulmonaria spores use some mechanism to inhibit germination—the inhibition is lifted when the spores are grown in a synthetic growth medium containing an adsorbent like bovine serum albumin or α-cyclodextrin.
Dispersal by vegetative propagules (via soredia or isidia) has been determined as the predominant mode of reproduction in L. pulmonaria. In this method, the protruding propagules become dry and brittle during the regular wet/dry cycles of the lichen, and can easily crumble off the thallus. These fragments may develop into new thalli, either at the same locale or at a new site after dispersal by wind or rain. A number of steps are required for the development of the vegetative propagules, including the degeneration of the thallus cortex, replication of green algal cells, and entanglement of fungal hyphae with the green algal cells. These steps lead to an increase in internal pressure which eventually breaks through the cortex. Continued growth leads to these granules being pushed upwards and out of the thallus surface.
Distribution and habitat
It has a wide distribution in Europe, Asia, North America and Africa, preferring damp habitats with high rainfall, especially coastal areas. It is the most widely distributed and most common Lobaria species in North America. In Wales, the Dolmelynllyn estate is notable for the variety of rare bryophytes and lichens there, including the genus Lobaria, in particular Lobaria pulmonaria. Associated with old-growth forests, its presence and abundance may be used as an indicator of forest age, at least in the Interior Cedar-Hemlock biogeoclimatic zone in eastern British Columbia. It is also found in pasture-woodlands. It usually grows on the bark of broad-leaved trees such as oak, beech and maple but will also grow on rocks. In the laboratory, L. pulmonaria has been grown on nylon microfilaments.
Various environmental factors are thought to affect the distribution of L. pulmonaria, such as temperature, moisture (average humidity, rapidity and frequency of wet-dry cycles), sunlight exposure, and levels of air pollution. Attempts to quantitatively evaluate the contribution of these factors to lichen growth is difficult because differences in the original environment from which the lichen thalli are collected will greatly affect heat and desiccation tolerances.
Due to declining population, L. pulmonaria is considered to be rare or threatened in many parts of the world, especially in lowland areas of Europe. The decline has been attributed to industrial forestry and air pollution, particularly acid rain. L. pulmonaria, like other lichens containing a blue-green algal component, are particularly susceptible to the effects of acid rain, because the subsequent decrease in pH reduces nitrogen fixation through inhibition of the algal nitrogenase enzyme.
Chemical compounds
Lobaria pulmonaria is known to contain a variety of acids common to lichens, such as stictic acid, desmethyl stictic acid, gyrophoric acid, tenuiorin, constictic acid, norstictic acid, peristictic acid, and methylnorstictic acid. These compounds, collectively known as depsidones, are known to be involved in defense against grazing herbivores like lichen-feeding molluscs. It also contains the sugar alcohols D-arabitol, volemitol, in addition to several carotenoids (total content > 10 mg/kg), such as alpha carotene, beta carotene, and beta cryptoxanthin. The upper cortex of the lichen contains melanins that screen UV and PAR radiation from the photobiont. The synthesis of melanin pigments in the lichen increases in response to greater solar irradiation, and shade-adapted thalli are greenish-grey in the air-dry state, while sun-exposed thalli can be dark brown in color. This adaptation helps protect the photosymbiont D. reticulata, known to be relatively intolerant to high light levels.
Also known to be present are various steroids, namely ergosterol, episterol, fecosterol, and lichesterol.
Uses
Medicinal
Its shape somewhat resembles the tissue inside lungs and therefore it is thought to be a remedy for lung diseases based on the doctrine of signatures. The lichen's common English names are derived from this association. Gerard's book The Herball or General Historie of plants (1597) recommends L. pulmonaria as medicinally valuable. It is still used for asthma, urinary incontinence and lack of appetite. In India it is used as a traditional medicine to treat hemorrhages and eczema, and it is used as a remedy for coughing up blood by the Hesquiaht in British Columbia, Canada. An ethnophytotherapeutical survey of the high Molise region in central-southern Italy revealed that L. pulmonaria is used as an antiseptic, and is rubbed on wounds.
A hot-water extract prepared using this species has been shown to have anti-inflammatory and ulcer-preventing activities. Also, methanol extracts were shown to have a protective effect on the gastrointestinal system of rats, possibly by reducing oxidative stress and reducing the inflammatory effects of neutrophils. Furthermore, methanol extracts also have potent antioxidative activity and reducing power, probably due to the presence of phenolic compounds.
Other uses
Lobaria pulmonaria has also been used to produce an orange dye for wool, in the tanning of leather, in the manufacture of perfumes and as an ingredient in brewing.
| Biology and health sciences | Lichens | Plants |
1827237 | https://en.wikipedia.org/wiki/Factorial%20experiment | Factorial experiment | In statistics, a full factorial experiment is an experiment whose design consists of two or more factors, each with discrete possible values or "levels", and whose experimental units take on all possible combinations of these levels across all such factors. A full factorial design may also be called a fully crossed design. Such an experiment allows the investigator to study the effect of each factor on the response variable, as well as the effects of interactions between factors on the response variable.
For the vast majority of factorial experiments, each factor has only two levels. For example, with two factors each taking two levels, a factorial experiment would have four treatment combinations in total, and is usually called a 2×2 factorial design. In such a design, the interaction between the variables is often the most important. This applies even to scenarios where a main effect and an interaction are present.
If the number of combinations in a full factorial design is too high to be logistically feasible, a fractional factorial design may be done, in which some of the possible combinations (usually at least half) are omitted.
Other terms for "treatment combinations" are often used, such as runs (of an experiment), points (viewing the combinations as vertices of a graph, and cells (arising as intersections of rows and columns).
History
Factorial designs were used in the 19th century by John Bennet Lawes and Joseph Henry Gilbert of the Rothamsted Experimental Station.
Ronald Fisher argued in 1926 that "complex" designs (such as factorial designs) were more efficient than studying one factor at a time. Fisher wrote,
A factorial design allows the effect of several factors and even interactions between them to be determined with the same number of trials as are necessary to determine any one of the effects by itself with the same degree of accuracy.
Frank Yates made significant contributions, particularly in the analysis of designs, by the Yates analysis.
The term "factorial" may not have been used in print before 1935, when Fisher used it in his book The Design of Experiments.
Advantages and disadvantages of factorial experiments
Many people examine the effect of only a single factor or variable. Compared to such one-factor-at-a-time (OFAT) experiments, factorial experiments offer several advantages
Factorial designs are more efficient than OFAT experiments. They provide more information at similar or lower cost. They can find optimal conditions faster than OFAT experiments.
When the effect of one factor is different for different levels of another factor, it cannot be detected by an OFAT experiment design. Factorial designs are required to detect such interactions. Use of OFAT when interactions are present can lead to serious misunderstanding of how the response changes with the factors.
Factorial designs allow the effects of a factor to be estimated at several levels of the other factors, yielding conclusions that are valid over a range of experimental conditions.
The main disadvantage of the full factorial design is its sample size requirement, which grows exponentially with the number of factors or inputs considered. Alternative strategies with improved computational efficiency include fractional factorial designs, Latin hypercube sampling, and quasi-random sampling techniques.
Example of advantages of factorial experiments
In his book, Improving Almost Anything: Ideas and Essays, statistician George Box gives many examples of the benefits of factorial experiments. Here is one. Engineers at the bearing manufacturer SKF wanted to know if changing to a less expensive "cage" design would affect bearing life. The engineers asked Christer Hellstrand, a statistician, for help in designing the experiment.
Box reports the following. "The results were assessed by an accelerated life test. … The runs were expensive because they needed to be made on an actual production line and the experimenters were planning to make four runs with the standard cage and four with the modified cage. Christer asked if there were other factors they would like to test. They said there were, but that making added runs would exceed their budget. Christer showed them how they could test two additional factors "for free" – without increasing the number of runs and without reducing the accuracy of their estimate of the cage effect. In this arrangement, called a 2×2×2 factorial design, each of the three factors would be run at two levels and all the eight possible combinations included. The various combinations can conveniently be shown as the vertices of a cube ... "
"In each case, the standard condition is indicated by a minus sign and the modified condition by a plus sign. The factors changed were heat treatment, outer ring osculation, and cage design. The numbers show the relative lengths of lives of the bearings. If you look at [the cube plot], you can see that the choice of cage design did not make a lot of difference. … But, if you average the pairs of numbers for cage design, you get the [table below], which shows what the two other factors did. … It led to the extraordinary discovery that, in this particular application, the life of a bearing can be increased fivefold if the two factor(s) outer ring osculation and inner ring heat treatments are increased together."
"Remembering that bearings like this one have been made for decades, it is at first surprising that it could take so long to discover so important an improvement. A likely explanation is that, because most engineers have, until recently, employed only one factor at a time experimentation, interaction effects have been missed."
Example
The simplest factorial experiment contains two levels for each of two factors. Suppose an engineer wishes to study the total power used by each of two different motors, A and B, running at each of two different speeds, 2000 or 3000 RPM. The factorial experiment would consist of four experimental units: motor A at 2000 RPM, motor B at 2000 RPM, motor A at 3000 RPM, and motor B at 3000 RPM. Each combination of a single level selected from every factor is present once.
This experiment is an example of a 22 (or 2×2) factorial experiment, so named because it considers two levels (the base) for each of two factors (the power or superscript), or #levels#factors, producing 22=4 factorial points.
Designs can involve many independent variables. As a further example, the effects of three input variables can be evaluated in eight experimental conditions shown as the corners of a cube.
This can be conducted with or without replication, depending on its intended purpose and available resources. It will provide the effects of the three independent variables on the dependent variable and possible interactions.
Notation
Factorial experiments are described by two things: the number of factors, and the number of levels of each factor. For example, a 2×3 factorial experiment has two factors, the first at 2 levels and the second at 3 levels. Such an experiment has 2×3=6 treatment combinations or cells. Similarly, a 2×2×3 experiment has three factors, two at 2 levels and one at 3, for a total of 12 treatment combinations. If every factor has s levels (a so-called fixed-level or symmetric design), the experiment is typically denoted by sk, where k is the number of factors. Thus a 25 experiment has 5 factors, each at 2 levels. Experiments that are not fixed-level are said to be mixed-level or asymmetric.
There are various traditions to denote the levels of each factor. If a factor already has natural units, then those are used. For example, a shrimp aquaculture experiment might have factors temperature at 25 °C and 35 °C, density at 80 or 160 shrimp/40 liters, and salinity at 10%, 25% and 40%. In many cases, though, the factor levels are simply categories, and the coding of levels is somewhat arbitrary. For example, the levels of an 6-level factor might simply be denoted 1, 2, ..., 6.
Treatment combinations are denoted by ordered pairs or, more generally, ordered tuples. In the aquaculture experiment, the ordered triple (25, 80, 10) represents the treatment combination having the lowest level of each factor. In a general 2×3 experiment the ordered pair (2, 1) would indicate the cell in which factor A is at level 2 and factor B at level 1. The parentheses are often dropped, as shown in the accompanying table.
To denote factor levels in 2k experiments, three particular systems appear in the literature:
The values 1 and 0;
the values 1 and −1, often simply abbreviated by + and −;
A lower-case letter with the exponent 0 or 1.
If these values represent "low" and "high" settings of a treatment, then it is natural to have 1 represent "high", whether using 0 and 1 or −1 and 1. This is illustrated in the accompanying table for a 2×2 experiment. If the factor levels are simply categories, the correspondence might be different; for example, it is natural to represent "control" and "experimental" conditions by coding "control" as 0 if using 0 and 1, and as 1 if using 1 and −1. An example of the latter is given below. That example illustrates another use of the coding +1 and −1.
For other fixed-level (sk) experiments, the values 0, 1, ..., s−1 are often used to denote factor levels. These are the values of the integers modulo s when s is prime.
Contrasts, main effects and interactions
The expected response to a given treatment combination is called a cell mean, usually denoted using the Greek letter μ. (The term cell is borrowed from its use in tables of data.) This notation is illustrated here for the 2 × 3 experiment.
A contrast in cell means is a linear combination of cell means in which the coefficients sum to 0. Contrasts are of interest in themselves, and are the building blocks by which main effects and interactions are defined.
In the 2 × 3 experiment illustrated here, the expression
is a contrast that compares the mean responses of the treatment combinations 11 and 12. (The coefficients here are 1 and –1.) The contrast
is said to belong to the main effect of factor A as it contrasts the responses to the "1" level of factor with those for the "2" level. The main effect of A is said to be absent if this expression equals 0.
Interaction in a factorial experiment is the lack of additivity between factors, and is also expressed by contrasts. In the 2 × 3 experiment, the contrasts
and
belong to the A × B interaction; interaction is absent (additivity is present) if these expressions equal 0. Additivity may be viewed as a kind of parallelism between factors, as illustrated in the Analysis section below.
Since it is the coefficients of these contrasts that carry the essential information, they are often displayed as column vectors. For the example above, such a table might look like this:
The columns of such a table are called contrast vectors: their components add up to 0. Each effect is determined by both the pattern of components in its columns and the number of columns.
The patterns of components of these columns reflect the general definitions given by Bose:
A contrast vector belongs to the main effect of a particular factor if the values of its components depend only on the level of that factor.
A contrast vector belongs to the interaction of two factors, say A and B, if (i) the values of its components depend only on the levels of A and B, and (ii) it is orthogonal (perpendicular) to the contrast vectors representing the main effects of A and B.
Similar definitions hold for interactions of more than two factors. In the 2 × 3 example, for instance, the pattern of the A column follows the pattern of the levels of factor A, indicated by the first component of each cell. Similarly, the pattern of the B columns follows the levels of factor B (sorting on B makes this easier to see).
The number of columns needed to specify each effect is the degrees of freedom for the effect, and is an essential quantity in the analysis of variance. The formula is as follows:
A main effect for a factor with s levels has s−1 degrees of freedom.
The interaction of two factors with s1 and s2 levels, respectively, has (s1−1)(s2−1) degrees of freedom.
The formula for more than two factors follows this pattern. In the 2 × 3 example above, the degrees of freedom for the two main effects and the interaction — the number of columns for each — are 1, 2 and 2, respectively.
Examples
In the tables in the following examples, the entries in the "cell" column are treatment combinations: The first component of each combination is the level of factor A, the second for factor B, and the third (in the 2 × 2 × 2 example) the level of factor C. The entries in each of the other columns sum to 0, so that each column is a contrast vector.
A 3 × 3 experiment: Here we expect 3-1 = 2 degrees of freedom each for the main effects of factors A and B, and (3-1)(3-1) = 4 degrees of freedom for the A × B interaction. This accounts for the number of columns for each effect in the accompanying table.
The two contrast vectors for A depend only on the level of factor A. This can be seen by noting that the pattern of entries in each A column is the same as the pattern of the first component of "cell". (If necessary, sorting the table on A will show this.) Thus these two vectors belong to the main effect of A. Similarly, the two contrast vectors for B depend only on the level of factor B, namely the second component of "cell", so they belong to the main effect of B.
The last four column vectors belong to the A × B interaction, as their entries depend on the values of both factors, and as all four columns are orthogonal to the columns for A and B. The latter can be verified by taking dot products.
A 2 × 2 × 2 experiment: This will have 1 degree of freedom for every main effect and interaction. For example, a two-factor interaction will have (2-1)(2-1) = 1 degree of freedom. Thus just a single column is needed to specify each of the seven effects.
The columns for A, B and C represent the corresponding main effects, as the entries in each column depend only on the level of the corresponding factor. For example, the entries in the B column follow the same pattern as the middle component of "cell", as can be seen by sorting on B.
The columns for AB, AC and BC represent the corresponding two-factor interactions. For example, (i) the entries in the BC column depend on the second and third (B and C) components of cell, and are independent of the first (A) component, as can be seen by sorting on BC; and (ii) the BC column is orthogonal to columns B and C, as can be verified by computing dot products.
Finally, the ABC column represents the three-factor interaction: its entries depend on the levels of all three factors, and it is orthogonal to the other six contrast vectors.
Combined and read row-by-row, columns A, B, C give an alternate notation, mentioned above, for the treatment combinations (cells) in this experiment: cell 000 corresponds to +++, 001 to ++−, etc.
In columns A through ABC, the number 1 may be replaced by any constant, because the resulting columns will still be contrast vectors.
For example, it is common to use the number 1/4 in 2 × 2 × 2 experiments to define each main effect or interaction, and to declare, for example, that the contrast
is "the" main effect of factor A, a numerical quantity that can be estimated.
Implementation
For more than two factors, a 2k factorial experiment can usually be recursively designed from a 2k−1 factorial experiment by replicating the 2k−1 experiment, assigning the first replicate to the first (or low) level of the new factor, and the second replicate to the second (or high) level. This framework can be generalized to, e.g., designing three replicates for three level factors, etc.
A factorial experiment allows for estimation of experimental error in two ways. The experiment can be replicated, or the sparsity-of-effects principle can often be exploited. Replication is more common for small experiments and is a very reliable way of assessing experimental error. When the number of factors is large (typically more than about 5 factors, but this does vary by application), replication of the design can become operationally difficult. In these cases, it is common to only run a single replicate of the design, and to assume that factor interactions of more than a certain order (say, between three or more factors) are negligible. Under this assumption, estimates of such high order interactions are estimates of an exact zero, thus really an estimate of experimental error.
When there are many factors, many experimental runs will be necessary, even without replication. For example, experimenting with 10 factors at two levels each produces 210=1024 combinations. At some point this becomes infeasible due to high cost or insufficient resources. In this case, fractional factorial designs may be used.
As with any statistical experiment, the experimental runs in a factorial experiment should be randomized to reduce the impact that bias could have on the experimental results. In practice, this can be a large operational challenge.
Factorial experiments can be used when there are more than two levels of each factor. However, the number of experimental runs required for three-level (or more) factorial designs will be considerably greater than for their two-level counterparts. Factorial designs are therefore less attractive if a researcher wishes to consider more than two levels.
Analysis
A factorial experiment can be analyzed using ANOVA or regression analysis. To compute the main effect of a factor "A" in a 2-level experiment, subtract the average response of all experimental runs for which A was at its low (or first) level from the average response of all experimental runs for which A was at its high (or second) level.
Other useful exploratory analysis tools for factorial experiments include main effects plots, interaction plots, Pareto plots, and a normal probability plot of the estimated effects.
When the factors are continuous, two-level factorial designs assume that the effects are linear. If a quadratic effect is expected for a factor, a more complicated experiment should be used, such as a central composite design. Optimization of factors that could have quadratic effects is the primary goal of response surface methodology.
Analysis example
Montgomery gives the following example of analysis of a factorial experiment:.An engineer would like to increase the filtration rate (output) of a process to produce a chemical, and to reduce the amount of formaldehyde used in the process. Previous attempts to reduce the formaldehyde have lowered the filtration rate. The current filtration rate is 75 gallons per hour. Four factors are considered: temperature (A), pressure (B), formaldehyde concentration (C), and stirring rate (D). Each of the four factors will be tested at two levels.Onwards, the minus (−) and plus (+) signs will indicate whether the factor is run at a low or high level, respectively.
The non-parallel lines in the A:C interaction plot indicate that the effect of factor A depends on the level of factor C. A similar results holds for the A:D interaction. The graphs indicate that factor B has little effect on filtration rate. The analysis of variance (ANOVA) including all 4 factors and all possible interaction terms between them yields the coefficient estimates shown in the table below.
Because there are 16 observations and 16 coefficients (intercept, main effects, and interactions), p-values cannot be calculated for this model. The coefficient values and the graphs suggest that the important factors are A, C, and D, and the interaction terms A:C and A:D.
The coefficients for A, C, and D are all positive in the ANOVA, which would suggest running the process with all three variables set to the high value. However, the main effect of each variable is the average over the levels of the other variables. The A:C interaction plot above shows that the effect of factor A depends on the level of factor C, and vice versa. Factor A (temperature) has very little effect on filtration rate when factor C is at the + level. But Factor A has a large effect on filtration rate when factor C (formaldehyde) is at the − level. The combination of A at the + level and C at the − level gives the highest filtration rate. This observation indicates how one-factor-at-a-time analyses can miss important interactions. Only by varying both factors A and C at the same time could the engineer discover that the effect of factor A depends on the level of factor C.
The best filtration rate is seen when A and D are at the high level, and C is at the low level. This result also satisfies the objective of reducing formaldehyde (factor C). Because B does not appear to be important, it can be dropped from the model. Performing the ANOVA using factors A, C, and D, and the interaction terms A:C and A:D, gives the result shown in the following table, in which all the terms are significant (p-value < 0.05).
| Mathematics | Statistics | null |
1827682 | https://en.wikipedia.org/wiki/Male%20reproductive%20system | Male reproductive system | The male reproductive system consists of a number of sex organs that play a role in the process of human reproduction. These organs are located on the outside of the body, and within the pelvis.
The main male sex organs are the penis and the scrotum, which contains the testicles that produce semen and sperm, which, as part of sexual intercourse, fertilize an ovum in the female's body; the fertilized ovum (zygote) develops into a fetus, which is later born as an infant. The corresponding system in females is the female reproductive system.
External genitalia
Penis
The penis is an intromittent organ with a long shaft, an enlarged bulbous-shaped tip called the glans and its foreskin for protection. Inside the penis is the urethra, which is used to ejaculate semen and to excrete urine. Both substances exit through the meatus.
When a male becomes sexually aroused, erection occurs because sinuses within the erectile tissues of the penis (corpora cavernosa and corpus spongiosum) become filled with blood. The arteries of the penis are dilated while the veins are compressed so that blood flows into the erectile cartilage under pressure. The penis is supplied by the pudendal artery.
Scrotum
The scrotum is a sac of skin that hangs behind the penis. It holds and protects the testicles. It also contains numerous nerves and blood vessels. During times of lower temperatures, the cremaster muscle contracts and pulls the scrotum closer to the body, while the dartos fascia gives it a wrinkled appearance; when the temperature increases, the cremaster and dartos fascia relax to bring down the scrotum away from the body and remove the wrinkles respectively.
The scrotum remains connected with the abdomen or pelvic cavity through the inguinal canal. (The spermatic cord, formed from spermatic artery, vein and nerve bound together with connective tissue passes into the testis through inguinal canal.)
Internal genitalia
Testicles
The testicles are two gonads that produce sperm by meiotic division of germ cells within the seminiferous tubules, and synthesize and secrete androgens that regulate the male reproductive functions. The site of production of androgens is the Leydig cells that are located in the interstitium between seminiferous tubules.
Epididymides
The epididymis is a long whitish mass of tightly coiled tube. The sperm that are produced in the seminiferous tubules flow into the epididymis. During passage via the epididymis, the sperm undergo maturation and are concentrated by the action of ion channels located on the apical membrane of the epididymis.
Vasa deferentia
The vas deferens, which is also known as the sperm duct, is a thin tube approximately long that starts from the epididymis to the pelvic cavity. It carries the spermatozoa from the epididymis to the ejaculatory duct.
Accessory glands
Three accessory glands provide fluids that lubricate the duct system and nourish the sperm cells.
Seminal vesicles: two glands behind the bladder that secrete many of the semen's components.
Prostate gland: a gland located below the bladder that produces seminal fluid and helps regulate urine flow.
Bulbourethral glands: add fluid to semen during ejaculation (pre-ejaculate).
Development
The embryonic and prenatal development of the male reproductive system is the process whereby the reproductive organs grow, mature and are established. It begins with a single fertilized egg and culminates 38 weeks later with the birth of a male child. It is a part of the stages of sexual differentiation. The development of the male reproductive system coincides with the urinary system. Their development can also be described together as the development of the urinary and reproductive organs.
Sexual determination
Sexual identity is determined at fertilization when the genetic sex of the zygote has been initialized by a sperm cell containing either an X or Y chromosome. If this sperm cell contains an X chromosome it will coincide with the X chromosome of the ovum and a female child will develop. A sperm cell carrying a Y chromosome results in an XY combination, and a male child will develop.
Genetic sex determines whether the gonads will be testes or ovaries. In the developing embryo if the testes are developed, it will produce and secrete male sex hormones during late embryonic development and cause the secondary sex organs of the male to develop.
Other embryonic reproductive structures
The structures are masculinized by secretions of the testes:
urogenital sinus
genital tubercle
urogenital folds
cloacal membrane
labioscrotal folds
The prostate gland derives from the urogenital sinus, and the other embryonic structures differentiate into the external genitalia. In the absence of testicular secretions, the female genitalia are formed.
External structures
At six weeks post-conception, the differentiation of the external genitalia in the male and female has not taken place. At eight weeks, a distinct phallus is present during the indifferent stage. By the 10th-12th week, the genitalia are distinctly male or female being and derived from their homologous structures. At 16 weeks post-conception, the genitalia are formed and distinct.
The masculinization of the embryonic reproductive structures occurs as a result of testosterone secreted by the embryonic testes. Testosterone, however, is not the active agent within these organs. Once inside the target cells, testosterone is converted by means of an enzyme called 5α-reductase into the dihydrotestosterone (DHT). DHT mediates the androgen effect in these organs.
Testes
At nine weeks, male differentiation of the gonads and the testes is well underway. Internal changes include the formation of the tubular seminar Chris tubules in the rete testis from the primary sex cord. Developing on the outside surface of each testis is a Phibro muscular cord called the gubernaculum. This structure attaches to the inferior portion of the testis and extends to the labial sacral fold of the same side at the same time, a portion of the embryonic mesonephric duct adjacent to the testis becomes attached and convoluted informs the epididymis. Another portion of the mesonephric duct becomes the ductus deferens.
The seminal vesicles form from lateral outgrowths of the caudal and of each mesonephric duct the prostate gland arises from an Indo dermal outgrowth of the urogenital sinus the bulbourethral glands develop from outgrowths in the membrane-like portion of the urethra.
The descent of the testes to its final location at the anterior abdominal wall, followed by the development of the gubernaculum, which subsequently pulls and translocates the testis down into the developing scrotum. Ultimately, the passageway closes behind the testis. A failure in this process can cause indirect inguinal hernia or an infantile hydrocoele. The testes descend into the scrotal sac between the sixth and 10th week. Descent does not occur until about the 28th week when compared to when canals form and the abdominal wall provides openings from the pelvic cavity to the scrotal sac. The process by which a testis descends is not well understood but it seems to be associated with the shortening of the gubernaculum. This is attached to the testis and extends through the inguinal canal to the wall of the scrotum as a testis. It carries with it the ductus deference, which are testicular vessels and nerves, a portion of the abdominal muscle, and lymph vessels. All of the structures remain attached to the testis and form what is known as the spermatic cord. By the time the testis is in the scrotal sac, the gubernaculum is no more than a remnant of scar like tissue.
Male germ cells formed in the testes are capable of special DNA repair processes that function during meiosis to repair DNA damages and to maintain the integrity of the genomes that are to be passed on to progeny. These DNA repair processes include homologous recombinational repair and non-homologous end joining.
External genitalia
The external genitalia of the male is distinct from those of the female by the end of the ninth week. Prior to that, the genital tubercle in both sexes is a phallus. The urethral groove forms on the ventral surface of the phallus early in development during the differentiation of the external genitalia. This is caused by the androgens produced and secreted by the testes. Androgen induced development causes the elongation and differentiation of the phallus into a penis, a fusion of the urogenital folds surrounding the urethral groove along the ventral surface of the penis, and a midline closure of the labioscrotal folds. This closure forms the wall of the scrotum the external genitalia. The external genitalia are completely formed by the end of the 12th week.
At birth, the development of the prepubertal male reproductive system is completed. During the second trimester of pregnancy, testosterone secretion in the male declines so that at birth the testes are inactive. Gonadotropin secretion is low until the beginning of puberty.
Summary
The genetic sex is determined by whether a Y bearing or next bearing sperm fertilizes the open; the presence or absence of a Y chromosome in turn determines whether the gonads of the embryo will be testes or ovaries; and the presence or absence of testes, finally, determines whether the sex accessory organs and external genitalia will be male or female. This sequence is understandable in light of the fact that both male and female embryos develop within the maternal environment - high in estrogen secreted by the mother's ovaries and the placenta. If estrogen determined the gender, all embryos would become feminized.
Puberty
During puberty, increased gonadotropin secretion stimulates a rise in sex steroids creation from the testes. The increased secretion of testosterone from the testes during puberty causes the male secondary sexual characteristics to be manifested.
Male secondary sex characteristics include:
Growth of body hair, including underarm, abdominal, chest hair and pubic hair.
Growth of facial hair.
Enlargement of larynx (Adam's apple) and deepening of voice.
Increased stature; adult males are taller than adult females, on average.
Heavier skull and bone structure.
Increased muscle mass and strength.
Broadening of shoulders and chest; shoulders wider than hips.
Increased secretions of oil and sweat glands.
Secondary development includes the increased activity of the eccrine sweat glands and sebaceous glands along with the darkening of the skin in the scrotal region.
Clinical significance
Chromosomal abnormalities
Chromosomal abnormalities can occur during fertilization impacting the development of the male reproductive system. The genotype of the male consists of a Y chromosome paired with an X chromosome. Female sex is determined by the absence of a Y chromosome. Some individuals are male who have the XX male syndrome and androgen insensitivity syndrome. This occurs when one X chromosome contains a segment of the Y chromosome, which was inserted into the X chromosome of the father's sperm. Rarely females are born with the XY genotype. They are found to be missing the same portion of the Y chromosome as was inserted into the chromosome of XX males. The gene for sexual differentiation in humans, called the testis determining factor (TDF), is located on the short arm of the Y chromosome. The presence or absence of the Y chromosome determines whether the embryo will have testes or ovaries. An abnormal number of sex chromosomes (aneuploidy) can occur. This includes Turner's syndrome - a single X chromosome is present, Klinefelter's syndrome - two X chromosomes and a Y chromosome are present, XYY syndrome and XXYY syndrome. Other less common chromosomal arrangements include: triple X syndrome, 48, XXXX, and 49, XXXXX.
The observable, visual differences become apparent between male or the female reproductive organs are not seen initially. Maturation continues as the medial aspect of each mesonephros grows to form the genital ridge. The genital ridge continues to grow behind the developing peritoneal membrane. By week six, string-like cell congregations called primitive sex cords form within the enlarging genital ridge. Externally, a swelling called the genital tubercle appears above the cloacal membrane.
External distinctions are not observed even by the eighth week of pre-embryonic development. This is the indifferent stage during which the gonads are relatively large and have an outer cortex of primitive sex cords and an inner medulla.
Specialized primordial germ cells are forming and migrating from the yolk sac to the embryonic gonads during week eight and nine. These are the spermatogonia in the developing male. Before seven weeks after fertilization, the gonads have the potential to become either testes or ovaries. Reproductive sex organs for both male and female are derived from the same embryonic tissues and are considered homologous tissues or organs.
After the testes have differentiated, male sex hormones, called androgens, are secreted from interstitial cells (cells of Leydig). The major androgens secreted by these cells is testosterone and secretion begins 8 to 10 weeks after conception. Testosterone secretion reaches a peak at 12 to 14 weeks, and declines to very low levels by the end of the second trimester (about 21 weeks). Levels are the barely detectable 4–6 months of age postnatal. High levels of testosterone will not appear again until the time of puberty.
Internal accessory sex organs to develop and most of these are derived from two systems of embryonic ducts. Male accessory organs are derived from mesonephric (wolfian) ducts. The developing tubules within the testes secretes a polypeptide Müllerian inhibition factor (MIF). MIF causes the regression of the paramesonephritic ducts 60 days after fertilization. Testosterone secretion by the interstitial cells of the testes then causes the growth and development of the mesonephric ducts into male secondary sex organs. The Müllerian ducts atrophy, but traces of their anterior ends are represented by the appendices testis (hydatids of Morgagni of the male), while their terminal fused portions form the utriculus on the floor of the prostatic urethra. This is due to the production of anti-Müllerian hormone by the Sertoli cells of the testes.
Gallery
| Biology and health sciences | Reproductive system | null |
1828131 | https://en.wikipedia.org/wiki/Davisson%E2%80%93Germer%20experiment | Davisson–Germer experiment | The Davisson–Germer experiment was a 1923–1927 experiment by Clinton Davisson and Lester Germer at Western Electric (later Bell Labs), in which electrons, scattered by the surface of a crystal of nickel metal, displayed a diffraction pattern. This confirmed the hypothesis, advanced by Louis de Broglie in 1924, of wave-particle duality, and also the wave mechanics approach of the Schrödinger equation. It was an experimental milestone in the creation of quantum mechanics.
History and overview
According to Maxwell's equations in the late 19th century, light was thought to consist of waves of electromagnetic fields and matter was thought to consist of localized particles. However, this was challenged in Albert Einstein's 1905 paper on the photoelectric effect, which described light as discrete and localized quanta of energy (now called photons), which won him the Nobel Prize in Physics in 1921. In 1924 Louis de Broglie presented his thesis concerning the wave–particle duality theory, which proposed the idea that all matter displays the wave–particle duality of photons. According to de Broglie, for all matter and for radiation alike, the energy of the particle was related to the frequency of its associated wave by the Planck relation:
And that the momentum of the particle was related to its wavelength by what is now known as the de Broglie relation:
where is the Planck constant.
An important contribution to the Davisson–Germer experiment was made by Walter M. Elsasser in Göttingen in the 1920s, who remarked that the wave-like nature of matter might be investigated by electron scattering experiments on crystalline solids, just as the wave-like nature of X-rays had been confirmed through Barkla's X-ray scattering experiments on crystalline solids.
This suggestion of Elsasser was then communicated by his senior colleague (and later Nobel Prize recipient) Max Born to physicists in England. When the Davisson and Germer experiment was performed, the results of the experiment were explained by Elsasser's proposition. However the initial intention of the Davisson and Germer experiment was not to confirm the de Broglie hypothesis, but rather to study the surface of nickel.
In 1927 at Bell Labs, Clinton Davisson and Lester Germer fired slow moving electrons at a crystalline nickel target. The angular dependence of the reflected electron intensity was measured and was determined to have a similar diffraction pattern as those predicted by Bragg for X-rays; some small, but significant differences were due to the average potential which Hans Bethe showed in his more complete analysis. At the same time George Paget Thomson and his student Alexander Reid independently demonstrated the same effect firing electrons through celluloid films to produce a diffraction pattern, and Davisson and Thomson shared the Nobel Prize in Physics in 1937. The exclusion of Germer from sharing the prize has puzzled physicists ever since. The Davisson–Germer experiment confirmed the de Broglie hypothesis that matter has wave-like behavior. This, in combination with the Compton effect discovered by Arthur Compton (who won the Nobel Prize for Physics in 1927), established the wave–particle duality hypothesis which was a fundamental step in quantum theory.
Early experiments
Davisson began work in 1921 to study electron bombardment and secondary electron emissions. A series of experiments continued through 1925.
Prior to 1923, Davisson had been working with Charles H. Kunsman on detecting the effects of electron bombardment on tungsten when they noticed that 1% of the electrons bounced straight back to the electron gun in elastic scattering. This small but unexpected result led Davisson to theorize that he could examine the electron configuration of the atom in an analogous manner to how the Rutherford alpha particle scattering had examined the nucleus. They changed to a high vacuum and used nickel along with various other metals with unimpressive results.
In October 1924 when Germer joined the experiment, Davisson’s actual objective was to study the surface of a piece of nickel by directing a beam of electrons at the surface and observing how many electrons bounced off at various angles. They expected that because of the small size of electrons, even the smoothest crystal surface would be too rough and thus the electron beam would experience diffused reflection.
The experiment consisted of firing an electron beam (from an electron gun, an electrostatic particle accelerator) at a nickel crystal, perpendicular to the surface of the crystal, and measuring how the number of reflected electrons varied as the angle between the detector and the nickel surface varied. The electron gun was a heated tungsten filament that released thermally excited electrons which were then accelerated through an electric potential difference, giving them a certain amount of kinetic energy, towards the nickel crystal. To avoid collisions of the electrons with other atoms on their way towards the surface, the experiment was conducted in a vacuum chamber. To measure the number of electrons that were scattered at different angles, a faraday cup electron detector that could be moved on an arc path about the crystal was used. The detector was designed to accept only elastically scattered electrons.
During the experiment, air accidentally entered the chamber, producing an oxide film on the nickel surface. To remove the oxide, Davisson and Germer heated the specimen in a high temperature oven, not knowing that this caused the formerly polycrystalline structure of the nickel to form large single crystal areas with crystal planes continuous over the width of the electron beam.
When they started the experiment again and the electrons hit the surface, they were scattered by nickel atoms in crystal planes (so the atoms were regularly spaced) of the crystal. This, in 1925, generated a diffraction pattern with unexpected and uncorrelated peaks due to the heating causing a ten crystal faceted area. They changed the experiment to a single crystal and started again.
Breakthrough
On his second honeymoon, Davisson attended the Oxford meeting of the British Association for the Advancement of Science in summer 1926. At this meeting, he learned of the recent advances in quantum mechanics. To Davisson's surprise, Max Born gave a lecture that used the uncorrelated diffraction curves from Davisson's 1923 research on platinum with Kunsman, using the data as confirmation of the de Broglie hypothesis of which Davisson was unaware.
Davisson then learned that in prior years, other scientists – Walter Elsasser, E. G. Dymond, and Blackett, James Chadwick, and Charles Ellis – had attempted similar diffraction experiments, but were unable to generate low enough vacuums or detect the low-intensity beams needed.
Returning to the United States, Davisson made modifications to the tube design and detector mounting, adding azimuth in addition to colatitude. Following experiments generated a strong signal peak at 65 V and an angle . He published a note to Nature titled, "The Scattering of Electrons by a Single Crystal of Nickel".
Questions still needed to be answered and experimentation continued through 1927, because Davisson was now familiar with the de Broglie formula and had designed the test to see if any effect could be discerned for a changed electron wavelength , according to the de Broglie relationship, which they knew should create a peak at 78 and not 65 V as their paper had shown. Because of their failure to correlate with the de Broglie formula, their paper introduced an ad hoc contraction factor of 0.7, which, however, could only explain eight of the thirteen beams.
By varying the applied voltage to the electron gun, the maximum intensity of electrons diffracted by the atomic surface was found at different angles. The highest intensity was observed at an angle with a voltage of 54 V, giving the electrons a kinetic energy of .
As Max von Laue proved in 1912, the periodic crystal structure serves as a type of three-dimensional diffraction grating. The angles of maximum reflection are given by Bragg's condition for constructive interference from an array, Bragg's law
for , , and for the spacing of the crystalline planes of nickel () obtained from previous X-ray scattering experiments on crystalline nickel.
According to the de Broglie relation, electrons with kinetic energy of have a wavelength of . The experimental outcome was via Bragg's law, which closely matched the predictions. As Davisson and Germer state in their 1928 follow-up paper to their Nobel prize winning paper, "These results, including the failure of the data to satisfy the Bragg formula, are in accord with those previously obtained in our experiments on electron diffraction. The reflection data fail to satisfy the Bragg relation for the same reason that the electron diffraction beams fail to coincide with their Laue beam analogues." However, they add, "The calculated wave-lengths are in excellent agreement with the theoretical values of h/mv as shown in the accompanying table." So although electron energy diffraction does not follow the Bragg law, it did confirm de Broglie's theory that particles behave like waves. The full explanation was provided by Hans Bethe who solved Schrödinger equation for the case of electron diffraction.
Davisson and Germer's accidental discovery of the diffraction of electrons was the first direct evidence confirming de Broglie's hypothesis that particles can have wave properties as well.
Davisson's attention to detail, his resources for conducting basic research, the expertise of colleagues, and luck all contributed to the experimental success.
Practical applications
The specific approach used by Davisson and Germer used low energy electrons, what is now called low-energy electron diffraction (LEED). It wasn't until much later that development of experimental methods exploiting ultra-high vacuum technologies (e.g. the approach described by Alpert in 1953) enabled the extensive use of LEED diffraction to explore the surfaces of crystallized elements and the spacing between atoms. Methods where higher energy electrons are used for diffraction in many different ways developed much earlier.
| Physical sciences | Quantum mechanics | Physics |
1828559 | https://en.wikipedia.org/wiki/Ornithomimus | Ornithomimus | Ornithomimus (; "bird mimic") is a genus of ornithomimid theropod dinosaurs from the Campanian and Maastrichtian ages of Late Cretaceous Western North America. Ornithomimus was a swift, bipedal dinosaur which fossil evidence indicates was covered in feathers and equipped with a small toothless beak that may indicate an omnivorous diet. It is usually classified into two species: the type species, Ornithomimus velox, and a referred species, Ornithomimus edmontonicus. O. velox was named in 1890 by Othniel Charles Marsh on the basis of a foot and partial hand from the Denver Formation of Colorado. Another seventeen species have been named since then, though almost all of them have been subsequently assigned to new genera or shown to be not directly related to Ornithomimus velox. The best material of species still considered part of the genus has been found in Alberta, representing the species O. edmontonicus, known from several skeletons from the Horseshoe Canyon Formation. Additional species and specimens from other formations are sometimes classified as Ornithomimus, such as Ornithomimus samueli (alternately classified in the genera Dromiceiomimus or Struthiomimus) from the earlier Dinosaur Park Formation.
History of discovery
First species named
The history of Ornithomimus classification and the classification of ornithomimids in general has been very complicated. The type species, Ornithomimus velox, was first named by O.C. Marsh in 1890 and is based on syntypes YPM 542 and YPM 548 (a partial hindlimb and forelimb, respectively), found by George Lyman Cannon in the Denver Formation of Colorado on June 30, 1889. The generic name means "bird mimic", derived from Greek words ὄρνις (ornis), "bird", and μῖμος (mimos), "mimic", in reference to the bird-like foot. The specific name means "swift" in Latin. Simultaneously, Marsh named two other species: Ornithomimus tenuis (based on specimen USNM 5814) and Ornithomimus grandis. Both consist of fragmentary fossils found by John Bell Hatcher in Montana, which is today understood as tyrannosauroid material. At first, Marsh assumed Ornithomimus was an ornithopod, but this changed when Hatcher found specimen USNM 4736, a partial ornithomimid skeleton, in Wyoming. Marsh named it Ornithomimus sedens in 1892. On that occasion, Ornithomimus minutus was also created based on specimen YPM 1049 (a metatarsus), but it has since been recognized as belonging to an alvarezsaurid.
A sixth species, Ornithomimus altus, was named in 1902 by Lawrence Lambe and was based on specimen CMN 930 (hindlimbs found in 1901 in Alberta), but this was renamed to a separate genus in 1916: Struthiomimus, by Henry Fairfield Osborn. In 1920, Charles Whitney Gilmore named Ornithomimus affinis for Dryosaurus grandis (Lull 1911), based on indeterminate material. In 1930, Loris Russell renamed Struthiomimus brevetertius (Parks 1926) and Struthiomimus samueli (Parks 1928) into Ornithomimus brevitertius and Ornithomimus samueli, respectively. The very same year, Oliver Perry Hay renamed Aublysodon mirandus (Leidy 1868) into Ornithomimus mirandus, which is today seen as a nomen dubium. In 1933, William Arthur Parks created the species Ornithomimus elegans, which is today seen as either Chirostenotes or Elmisaurus. That same year, Gilmore named Ornithomimus asiaticus for material found in Inner Mongolia.
Also in 1933, Charles Mortram Sternberg named the species Ornithomimus edmontonicus for a nearly complete skeleton from the Horseshoe Canyon Formation (specimen CMN 8632).
Reclassification by Dale Russell
At first, it had been common practice to name each newly discovered ornithomimid as a species of Ornithomimus. In the sixties, this tendency was still very strong, as is shown by the fact that Oskar Kuhn renamed Megalosaurus lonzeensis (Dollo 1903) from Belgium into Ornithomimus lonzeensis (which is understood today to be an abelisauroid claw) and Dale Russell in 1967 renamed Struthiomimus currellii (Parks 1933) and Struthiomimus ingens (Parks 1933) into Ornithomimus currellii and Ornithomimus ingens, respectively. At the same time, it was usual that workers referred to the entire ornithomimid material as simply "Struthiomimus". To solve this confusion by scientifically testing the separation between Ornithomimus and Struthiomimus, Dale Russell in 1972 published a morphometric study. It showed that statistical differences in some proportions could be used to distinguish the two and he concluded that Struthiomimus and Ornithomimus were valid genera. In the latter, Russell recognised two species: the type species Ornithomimus velox and Ornithomimus edmontonicus (even though he had trouble reliably distinguishing it from O. velox). He considered Struthiomimus currellii to be a younger synonym of Ornithomimus edmontonicus. However, Russell also interpreted the data as indicating that many specimens could not be referred to either Ornithomimus or Struthiomimus. Therefore, he created two new genera. The first one was Archaeornithomimus. Ornithomimus asiaticus and Ornithomimus affinis were reassigned to this new genus, becoming Archaeornithomimus asiaticus and Archaeornithomimus affinis. The second one was Dromiceiomimus, meaning "Emu mimic". This comes from the old generic name for the emu: Dromiceius. Russell assigned several former Ornithomimus species named during the 20th century, including O. brevitertius and O. ingens, to this new genus as Dromiceiomimus brevitertius. He also renamed Ornithomimus samueli into a second Dromiceiomimus species: Dromiceiomimus samueli.
Misassigned to Ornithomimus
Two tibiae from the Navesink Formation of New Jersey were named Coelosaurus antiquus ("antique hollow lizard") by Joseph Leidy in 1865. The tibiae were first attributed to Ornithomimus in 1979 by Donald Baird and John R. Horner as Ornithomimus antiquus. Normally, this would have made Ornithomimus a junior synonym of Coelosaurus, but Baird and Horner discovered that the name "Coelosaurus" was preoccupied by a dubious taxon, which was based on a single vertebra. It was originally named Coelosaurus by an anonymous author now known to be Richard Owen in 1854. Baird referred several other specimens from New Jersey and Maryland to O. antiquus. Beginning in 1997, Robert M. Sullivan regarded O. velox and O. edmontonicus as junior synonyms of O. antiquus. Like Russell, he considered the former two species indistinguishable from each other and noted that they both shared distinctive features with O. antiquus. However, David Weishampel (2004) considered "C." antiquus to be indeterminate among ornithomimosaurs, resulting in it being a nomen dubium. An SVP 2012 abstract agreed with Weishampel by noting that Coelosaurus differs from Gallimimus and Ornithomimus in the features of the tibiae.
In 1988, Gregory S. Paul classified the various species of Archaeornithomimus, Struthiomimus, Dromiceiomimus, and Gallimimus to the genus Ornithomimus. This has found no acceptance among other workers and the name is not presently used by Paul himself.
Present interpretations
Even after Russell's study, various researchers have found reasons to lump some or all of these species back into Ornithomimus in various combinations. In 2004, Peter Makovicky, Yoshitsugu Kobayashi, and Phil Currie studied Russell's 1972 proportional statistics to re-analyze ornithomimid relationships in light of newly discovered specimens. They concluded that there was no justification to separate Dromiceiomimus from Ornithomimus, sinking Dromiceiomimus as a synonym of O. edmontonicus. However, they did not include the type species O. velox in this analysis. The same team further supported the synonymy between Dromiceiomimus and O. edmontonicus in a 2006 lecture at the Society of Vertebrate Paleontology annual meeting. Their opinion has been followed by most later authors. Makovicky's team also considered Dromiceiomimus samueli to be a junior synonym of O. edmontonicus, though Longrich later suggested it may belong to a distinct, unnamed species from the Dinosaur Park Formation which has yet to be described. Longrich called the species Ornithomimus samueli in a faunal list for the Dinosaur Park Formation.
Apart from O. edmontonicus dating to the early Maastrichtian, two other species are presently considered to be possibly valid and are also from the late Maastrichtian. O. sedens was named by Marsh in 1892 from partial remains found in the Lance Formation of Wyoming a year after the description of O. velox. Dale Russell, in his 1972 revision of ornithomimids, could not determine which genus it actually belonged to, though he speculated that it may be intermediate between Struthiomimus and Dromiceiomimus. In 1985, he considered it to be a species of Ornithomimus. Although it has since been referred to mainly as Struthiomimus sedens (based on complete specimens from Montana and some fragments from Alberta and Saskatchewan), these have yet to be described and compared to the O. sedens holotype.
The other is the original type species: O. velox, at first known from very limited remains. Additional specimens referred to O. velox have been described from the Denver Formation and from the Ferris Formation of Wyoming. One specimen attributed to O. velox (MNA P1 1762A) from the Kaiparowits Formation of Utah was described in 1985. Re-evaluation of this specimen by Lindsay Zanno and colleagues in 2010, however, cast doubt on its assignment to O. velox and possibly even to Ornithomimus. This conclusion was supported by a 2015 re-description of O. velox, which found that only the holotype specimen was confidently referable to that species. The authors of this study tentatively referred to the Kaiparowits specimen as Ornithomimus sp., along with all of the specimens from the Dinosaur Park Formation.
Description
Like other ornithomimids, species of Ornithomimus are characterized by feet with three weight-bearing toes, long slender arms, and long necks with birdlike, elongated, toothless, beaked skulls. They were bipedal and superficially resemblant to ratites. They would have been swift runners thanks to their very long limbs and hollow bones. They also had large brains and large eyes. The brains of ornithomimids in general were large for non-avialan dinosaurs, but this may not necessarily be a sign of high intelligence. Some paleontologists think that the enlarged portions of the brain were dedicated to kinesthetic coordination. The bones of the hands are remarkably sloth-like in appearance, which led Henry Fairfield Osborn to suggest that they were used to hook branches during feeding.
Ornithomimus differs from other ornithomimids, such as Struthiomimus, in having shorter torsos, long slender forearms, very slender, straight hand and foot claws, and hand bones and fingers of similar lengths.
The two Ornithomimus species today seen as possibly valid differ in size. In 2010, Gregory S. Paul estimated the length of O. edmontonicus at and its weight at . One of its specimens, CMN 12228, preserves a femur that is long. O. velox, the type species of Ornithomimus, is based on material of a smaller animal. Whereas the holotype of O. edmontonicus, CMN 8632, preserves a second metacarpal eighty-four millimetres long, the same element with O. velox measures only fifty-three millimetres.
Feathers and skin
Ornithomimus, like many dinosaurs, was long thought to have been scaly. However, beginning in 1995, several specimens of Ornithomimus have been found preserving evidence of feathers. In 1995, 2008, and 2009, three Ornithomimus edmontonicus specimens with evidence of feathers were found (two adults with carbonized traces on the lower arm, indicating the former presence of pennaceous feather shafts, and a juvenile with impressions of feathers, which were up to five centimetres in length, in the form of hair-like filaments covering the rump, legs, and neck was also discovered). The fact that the feather imprints were found in sandstone, previously thought to not be able to support such impressions, raised the possibility of finding similar structures with more careful preparation of future specimens. A study describing the fossils in 2012 concluded that O. edmontonicus was covered in plumaceous feathers at all growth stages and that only adults had pennaceous wing-like structures, suggesting that wings may have evolved for mating displays. In 2014, Christian Foth and others argued that the evidence was insufficient to conclude that the forelimb feathers of Ornithomimus were necessarily pennaceous, citing the fact that the monofilamentous wing feathers in cassowaries would likely leave similar traces.
A fourth feathered specimen of Ornithomimus, this time from the lower portion of the Dinosaur Park Formation, was described in October of 2015 by Aaron van der Reest, Alex Wolfe, and Phil Currie. It was the first Ornithomimus specimen to preserve the feathers of its tail. The feathers, though crushed and distorted, bore numerous similarities to those of an ostrich, both in structure and distribution. Skin impressions were also preserved in the 2015 specimen, which indicated that, from mid-thigh to the feet, there was bare skin and that a flap of skin connecting the upper thigh to the torso. This latter structure is similar to that found in modern birds, including ostriches, but was positioned higher above the knee in Ornithomimus than in birds.
Classification
In 1890, Marsh assigned Ornithomimus to the clade Ornithomimosauria, a classification that is still very common. Modern cladistic studies indicate it having a derived position in the ornithomimids. These, however, have only included O. edmontonicus in their analyses. The relationships between O. edmontonicus, O. velox, and O. sedens have not been published.
The following cladogram is based on Xu et al., 2011:
Paleobiology
The diet of Ornithomimus is still highly debated. As theropods, ornithomimids might have been carnivorous, but their body shape would have been suited for a largely omnivorous lifestyle. Suggested food in its diet includes insects, crustaceans, fruit, leaves, branches, eggs, lizards, and small mammals.
Ornithomimus had legs that seem clearly suited for rapid locomotion, with the tibia being about 20% longer than the femur. The large eye sockets suggest a keen visual sense and also suggest the possibility that they were nocturnal.
In a 2001 study conducted by Bruce Rothschild and other paleontologists, 178 foot bones referred to Ornithomimus were examined for signs of stress fracture, but none were found.
| Biology and health sciences | Theropods | Animals |
1828723 | https://en.wikipedia.org/wiki/Neomycin/polymyxin%20B/bacitracin | Neomycin/polymyxin B/bacitracin | Neomycin/polymyxin B/bacitracin, also known as triple antibiotic ointment, is an antibiotic medication used to reduce the risk of infections following minor skin injuries. It contains the three antibiotics neomycin, polymyxin B, and bacitracin. It is for topical use.
Possible side effects include itchiness and skin rash, and in rare cases hearing loss. It is relatively broad spectrum, being effective against both Gram-negative and Gram-positive bacteria.
The combination is available over the counter in the US and Canada. In 2021, it was the 376th most commonly prescribed medication in the United States, with more than 25,000 prescriptions.
Medical uses
Neomycin/polymyxin B/bacitracin ointment is reported to be a safe and effective topical agent for preventing infections in minor skin trauma.
It is used for burns, scratches, cuts, and minor skin infections.
The use of neomycin/polymyxin B/bacitracin, decreases infection rates in minor-contaminated wounds.
It is for external use only.
Side effects
It has been shown to cause contact dermatitis in some cases.
Antibiotic-resistant bacteria
Concern exists that its use contributes to the emergence of antibiotic-resistant bacteria. In the US, the only large market for the ointment, it may increase antibiotic resistance. For instance, it may increase the prevalence of methicillin-resistant Staphylococcus aureus (MRSA) bacteria, specifically the highly lethal ST8:USA300 strain.
Components
The 2023 updated Johnson & Johnson Consumer Inc. label for their product discloses three different antibiotics: bacitracin zinc 400 units, neomycin sulfate 3.5 mg, and polymyxin B sulfate 5,000 units, in a relatively low-molecular-weight base of petroleum jelly, cottonseed oil, olive oil, and cocoa butter, and with sodium pyruvate and tocopheryl acetate.
The generic name for these products, regardless of the base, is "triple antibiotic ointment".
In China, this product (with lidocaine HCl) is named "FONOW® Ointment (孚诺®软膏, Compound Polymyxin B Ointment)" and is exclusively manufactured and sold by Zhejiang Fonow Medicine Co. Ltd. The product was also marketed by the Upjohn Company under the name "Mycitracin", until 1997 when that name was acquired by Johnson & Johnson.
Some people have allergic reactions to neomycin, so a "double antibiotic ointment" is sold without it, containing only
bacitracin and polymyxin B: one such example is Polysporin branded product.
A variant of Polysporin, called Polysporin Triple Ointment, replaces neomycin with gramicidin, providing an alternative for those allergic to neomycin while still offering broad-spectrum coverage against both Gram-negative and Gram-positive bacteria.
Active ingredients
The three main active ingredients in Neosporin are neomycin sulfate, polymyxin B sulfate, and bacitracin zinc.
One of the main components is neomycin sulfate, which is a type of antibiotic discovered in 1949 by microbiologist Selman Waksman at Rutgers University. Neomycin belongs to the aminoglycoside class of antibiotics and fights against Gram positive and gram negative bacteria. The antibiotic is often used to prevent risk of bacterial infections. Aminoglycosides work by binding to bacterial RNA and changing the ability to produce proteins while exerting little to no effect on DNA. Thus, neomycin kills bacteria as a result of irregular protein production in the bacterial cell. When the cell can no longer produce the correct proteins, its membrane becomes damaged. As a result of damaged membrane, the affected bacterial cells die, and the infection is prevented or limited.
Pramoxine is used to temporarily reduce pain from burns, insect bites, and minor cuts. It works like an anesthetic by decreasing the permeability of neuron membranes. As a result, pain neurons in the area have difficulty sending signals (or signals are blocked entirely), resulting in numbness.
In some countries bacitracin is replaced with gramicidin. The original Neosporin was using this combination.
History
There is no exact date as to when the antibacterial ointment was invented, but it was used as early as the 1950s. This antibiotic ointment was patented in the United States in August 1952.
The brand Neosporin was first used in commerce in August 1952, and trademarked in October 1952.
| Biology and health sciences | Antibiotics | Health |
1829229 | https://en.wikipedia.org/wiki/Padlock | Padlock | Padlocks are portable locks with a shackle that may be passed through an opening (such as a chain link, or hasp staple) to prevent use, theft, vandalism or harm.
Naming and etymology
The term padlock is from the late fifteenth century. The prefix pad- is thought to be related to the Latin which may refer to the portability of a padlock; it is combined with the noun lock, from Old English , related to German , "hole".
History
There are padlocks dating to the Roman Era, 500 BC – 300 AD. They were known in early times by merchants traveling the ancient trade routes to Asia, including China.
Padlocks have been used in Europe since the middle La Tène period, subsequently spreading to the Roman world and the Przeworsk and Chernyakhov cultures. Roman padlocks had a long bent rod attached to the case, and a shorter piece which could be inserted into the case. Przeworsk and Chernyakhov padlocks had a sleeve attached to the case, and a long bent rod which could be inserted into the case and the sleeve.
Padlocks have been used in China since the late Eastern Han dynasty (25–220 AD). According to Hong-Sen Yan, director of the National Science and Technology Museum, early Chinese padlocks were mainly "key-operated locks with splitting springs, and partially keyless letter combination locks". Padlocks were made from bronze, brass, silver, and other materials. The use of bronze was more prevalent for the early Chinese padlocks.
Padlocks with spring tine mechanisms have been found in York, England, at the Jorvik Viking settlement, dated 850 AD.
Smokehouse locks, designed in England, were formed from wrought iron sheet and employed simple lever and ward mechanisms. These locks afforded little protection against forced and surreptitious entry. Contemporary with the smokehouse padlocks and originating in the Slavic areas of Europe, "screw key" padlocks opened with a helical key that was threaded into the keyhole. The key pulled the locking bolt open against a strong spring. Padlocks that offered more key variance were the demise of the screw lock. Improved manufacturing methods allowed the manufacture of better padlocks that put an end to the Smokehouse around 1910.
Around the middle of the 19th century, "Scandinavian" style locks, or "Polhem locks", invented by the eponymous Swedish inventor Christopher Polhem, became a more secure alternative to the prevailing smokehouse and screw locks. These locks had a cast iron body that was loaded with a stack of rotating disks. Each disk had a central cutout to allow the key to pass through them and two notches cut out on the edge of the disc. When locked, the discs passed through cut-outs on the shackle. The key rotated each disk until the notches, placed along the edge of each tumbler in different places, lined up with the shackle, allowing the shackle to slide out of the body. The McWilliams company received a patent for these locks in 1871. The "Scandinavian" design was so successful that JHW Climax & Co. of Newark, New Jersey continued to make these padlocks until the 1950s. Today, other countries are still manufacturing this style of padlock.
Contemporary with the Scandinavian padlock, were the "cast heart" locks, so called because of their shape. A significantly stronger lock than the smokehouse and much more resistant to corrosion than the Scandinavian, the hearts had a lock body sand cast from brass or bronze and a more secure lever mechanism. Heart locks had two prominent characteristics: one was a spring-loaded cover that pivoted over the keyhole to keep dirt and insects out of the lock that was called a "drop". The other was a point formed at the bottom of the lock so a chain could be attached to the lock body to prevent the lock from getting lost or stolen. Cast heart locks were very popular with railroads for locking switches and cars because of their economical cost and excellent ability to open reliably in dirty, moist, and frozen environments.
Around the 1870s, lock makers realized they could successfully package the same locking mechanism found in cast heart locks into a more economical steel or brass shell instead of having to cast a thick metal body. These lock shells were stamped out of flat metal stock, filled with lever tumblers, and then riveted together. Although more fragile than the cast hearts, these locks were attractive because they cost less. In 1908, Adams & Westlake patented a stamped & riveted switch lock that was so economical that many railroads stopped using the popular cast hearts and went with this new stamped shell lock body design. Many lock manufacturers made this very popular style of lock.
In 1877 Yale & Towne was granted a patent for a padlock that housed a stack of levers and had a shackle that swung away when unlocked. It was a notable design because the levers were sub-assembled into a "cartridge" that could be slid into a cast brass body shell. The assembly would remain together by means of two taper pins passed through the shell and cartridge. This design gave the commercial padlock market a serviceable, rekeyable padlock. About twenty years later Yale made another "cartridge" style padlock that employed their famous pin tumbler mechanism and a shackle that slid out of the body instead of swinging away.
Although machining metal was a method that was available to lock makers since the early 19th century, it was not economically feasible to do so until the very early 20th century when electrical generation and distribution became widespread. Some of the earliest padlocks (c. 1905) that were made from a machined block of cast or extruded metal resemble today's modern padlock. Corbin and Eagle were one of the first lock makers to machine a solid block of metal and insert a relatively new pin tumbler mechanism and a sliding shackle into the holes machined into the body. This style of padlock was both strong and easy to manufacture. Many machined body padlocks were designed to be disassembled so that locksmiths could easily fit the locks to a certain key. The machined body padlocks are still very popular today. The process of machining allows many modern padlocks to have a "shroud" covering the shackle, which is an extension of the body around the shackle to protect the shackle from getting sheared or cut.
In the early 1920s, Harry Soref started Master Lock off with the first laminated padlock. Plates that were punched from sheet metal were stacked and assembled. Holes that were formed in the middle of the plates made room to accommodate the locking mechanism. The entire stack of plates, loaded with the lock parts in it, was riveted together. This padlock was popular for its low cost and impact-resistant laminated plate design. Today, many lock makers copy this very efficient and successful design.
Components
A padlock is composed of a body, shackle, and locking mechanism. The typical shackle is a U-shaped loop of metal (round or square in cross-section) that encompasses what is being secured by the padlock (e.g., chain link or hasp). Generally, most padlock shackles either swing away (typical of older padlocks) or slide out of the padlock body when in the unlocked position. Less common designs include a straight, circular, or flexible (cable) shackle. Some shackles split apart and come together to lock and unlock.
There are two basic types of padlock locking mechanisms: integrated and modular. Integrated locking mechanisms directly engage the padlock's shackle with the tumblers. Examples of integrated locking mechanisms are rotating disks (found in "Scandinavian" style padlocks where a disk rotated by the key enters a notch cut into the shackle to block it from moving) or lever tumblers (where a portion of the bolt that secures the shackle enters the tumblers when the correct key is turned in the lock). Padlocks with integrated locking mechanisms are characterized by a design that does not allow disassembly of the padlock. They are usually older than padlocks with modular mechanisms and often require the use of a key to lock.
The more modern modular locking mechanisms, however, do not directly employ the tumblers to lock the shackle. Instead, they have a plug within the "cylinder" that, with the correct key, turns and allows a mechanism, referred to as a "locking dog" (such as the ball bearings found in American Lock Company padlocks) to retract from notches cut into the shackle. Padlocks with modular locking mechanisms can often be taken apart to change the tumblers or to service the lock. Modular locking mechanism cylinders frequently employ pin, wafer, and disc tumblers. Padlocks with modular mechanisms are usually automatic, or self-locking (that is, the key is not required to lock the padlock)
Types
Combination
Combination locks do not use keys. Instead, the lock opens when its wheels are lined up correctly to display the correct combination.
A padlock was invented by John I. W. Carlson in 1931 (a patent was granted in November 1934) that has both a combination on one side and a key on the other.
Electronic
Electronic padlocks are categorized as either "active" or "passive" electronic padlocks.
Active electronic padlocks are divided into rechargeable electronic padlocks and replaceable battery electronic padlocks. The key of this type of electronic padlock is no longer a key in the traditional sense, but is unlocked through mobile phone Bluetooth, near-field communication (NFC), and fingerprint.
Unlike active electronic padlocks, passive electronic padlocks do not require electricity which makes their use environment more extensive. The passive electronic lock can be unlocked with an electronic key which needs to be programmed.
Steel cable
Padlocks with a flexible cable shackle have been made to encompass larger or irregular things such as a bicycle, gate, or even the zippers on backpacks and luggage. Specialized cable shackle padlocks are made with a high security cylinder to secure keys to prevent keys from being removed. Typically used by vending companies, such cable shackle padlocks prevent the surreptitious removal of a key for duplication or unauthorized use outside of working hours.
Resistance
A quantitative measure of a padlock's tensile strength and resistance to forced and surreptitious entry can be determined with tests developed by organizations such as ASTM, Sold Secure (United Kingdom), CEN (Europe), and TNO (The Netherlands).
Symbolism
Images of closed padlocks (sometimes physical padlocks) are sometimes used to indicate that something is secure or inaccessible.
A widely known example is on the Web, where data in secure transactions is encrypted using public-key cryptography; some web browsers display a locked padlock icon while using the HTTPS protocol.
Love locks are physical padlocks attached to fixtures such as bridges, gates, and monuments by sweethearts to declare their love for each other is permanent. This practice perhaps originated in Serbia. In some popular tourist locations, huge numbers of such padlocks are treated as a nuisance by local authorities.
| Technology | Mechanisms | null |
1830142 | https://en.wikipedia.org/wiki/Bijection%2C%20injection%20and%20surjection | Bijection, injection and surjection | In mathematics, injections, surjections, and bijections are classes of functions distinguished by the manner in which arguments (input expressions from the domain) and images (output expressions from the codomain) are related or mapped to each other.
A function maps elements from its domain to elements in its codomain. Given a function :
The function is injective, or one-to-one, if each element of the codomain is mapped to by at most one element of the domain, or equivalently, if distinct elements of the domain map to distinct elements in the codomain. An injective function is also called an injection. Notationally:
or, equivalently (using logical transposition),
The function is surjective, or onto, if each element of the codomain is mapped to by at least one element of the domain; that is, if the image and the codomain of the function are equal. A surjective function is a surjection. Notationally:
The function is bijective (one-to-one and onto, one-to-one correspondence, or invertible) if each element of the codomain is mapped to by exactly one element of the domain; that is, if the function is both injective and surjective. A bijective function is also called a bijection. That is, combining the definitions of injective and surjective,
where means "there exists exactly one ".
In any case (for any function), the following holds:
An injective function need not be surjective (not all elements of the codomain may be associated with arguments), and a surjective function need not be injective (some images may be associated with more than one argument). The four possible combinations of injective and surjective features are illustrated in the adjacent diagrams.
Injection
A function is injective (one-to-one) if each possible element of the codomain is mapped to by at most one argument. Equivalently, a function is injective if it maps distinct arguments to distinct images. An injective function is an injection. The formal definition is the following.
The function is injective, if for all ,
The following are some facts related to injections:
A function is injective if and only if is empty or is left-invertible; that is, there is a function such that identity function on X. Here, is the image of .
Since every function is surjective when its codomain is restricted to its image, every injection induces a bijection onto its image. More precisely, every injection can be factored as a bijection followed by an inclusion as follows. Let be with codomain restricted to its image, and let be the inclusion map from into . Then . A dual factorization is given for surjections below.
The composition of two injections is again an injection, but if is injective, then it can only be concluded that is injective (see figure).
Every embedding is injective.
Surjection
A function is surjective or onto if each element of the codomain is mapped to by at least one element of the domain. In other words, each element of the codomain has a non-empty preimage. Equivalently, a function is surjective if its image is equal to its codomain. A surjective function is a surjection. The formal definition is the following.
The function is surjective, if for all , there is such that
The following are some facts related to surjections:
A function is surjective if and only if it is right-invertible; that is, if and only if there is a function such that identity function on . (This statement is equivalent to the axiom of choice.)
By collapsing all arguments mapping to a given fixed image, every surjection induces a bijection from a quotient set of its domain to its codomain. More precisely, the preimages under f of the elements of the image of are the equivalence classes of an equivalence relation on the domain of , such that x and y are equivalent if and only they have the same image under . As all elements of any one of these equivalence classes are mapped by on the same element of the codomain, this induces a bijection between the quotient set by this equivalence relation (the set of the equivalence classes) and the image of (which is its codomain when is surjective). Moreover, f is the composition of the canonical projection from f to the quotient set, and the bijection between the quotient set and the codomain of .
The composition of two surjections is again a surjection, but if is surjective, then it can only be concluded that is surjective (see figure).
Bijection
A function is bijective if it is both injective and surjective. A bijective function is also called a bijection or a one-to-one correspondence (not to be confused with one-to-one function, which refers to injection). A function is bijective if and only if every possible image is mapped to by exactly one argument. This equivalent condition is formally expressed as follows:
The function is bijective, if for all , there is a unique such that
The following are some facts related to bijections:
A function is bijective if and only if it is invertible; that is, there is a function such that identity function on and identity function on . This function maps each image to its unique preimage.
The composition of two bijections is again a bijection, but if is a bijection, then it can only be concluded that is injective and is surjective (see the figure at right and the remarks above regarding injections and surjections).
The bijections from a set to itself form a group under composition, called the symmetric group.
Cardinality
Suppose that one wants to define what it means for two sets to "have the same number of elements". One way to do this is to say that two sets "have the same number of elements", if and only if all the elements of one set can be paired with the elements of the other, in such a way that each element is paired with exactly one element. Accordingly, one can define two sets to "have the same number of elements"—if there is a bijection between them. In which case, the two sets are said to have the same cardinality.
Likewise, one can say that set "has fewer than or the same number of elements" as set , if there is an injection from to ; one can also say that set "has fewer than the number of elements" in set , if there is an injection from to , but not a bijection between and .
Examples
It is important to specify the domain and codomain of each function, since by changing these, functions which appear to be the same may have different properties.
Injective and surjective (bijective)
The identity function idX for every non-empty set X, and thus specifically
, and thus also its inverse
The exponential function (that is, the exponential function with its codomain restricted to its image), and thus also its inverse the natural logarithm Here, denotes the positive real numbers.
Injective and non-surjective
The exponential function
Non-injective and surjective
Non-injective and non-surjective
Properties
For every function , let be a subset of the domain and a subset of the codomain. One has always and , where is the image of and is the preimage of under . If is injective, then , and if is surjective, then .
For every function , one can define a surjection and an injection . It follows that . This decomposition as the composition of a surjection and an injection is unique up to an isomorphism, in the sense that, given such a decomposition, there is a unique bijection such that and for every
Category theory
In the category of sets, injections, surjections, and bijections correspond precisely to monomorphisms, epimorphisms, and isomorphisms, respectively.
History
The Oxford English Dictionary records the use of the word injection as a noun by S. Mac Lane in Bulletin of the American Mathematical Society (1950), and injective as an adjective by Eilenberg and Steenrod in Foundations of Algebraic Topology (1952).
However, it was not until the French Bourbaki group coined the injective-surjective-bijective terminology (both as nouns and adjectives) that they achieved widespread adoption.
| Mathematics | Functions: General | null |
4685590 | https://en.wikipedia.org/wiki/Optical%20path | Optical path | Optical path (OP) is the trajectory that a light ray follows as it propagates through an optical medium.
The geometrical optical-path length or simply geometrical path length (GPD) is the length of a segment in a given OP, i.e., the Euclidean distance integrated along a ray between any two points.
The mechanical length of an optical device can be reduced to less than the GPD by using folded optics.
The optical path length in a homogeneous medium is the GPD multiplied by the refractive index of the medium.
Factors affecting optical path
Path of light in medium, or between two media is affected by the following:
Reflection
Total internal reflection
Refraction
Dispersion of light
Absorption
Simple materials used
Lenses
Prisms
Mirrors
Transparent materials (e.g. optical filters)
Translucent materials (e.g. frosted glass)
Opaque materials
| Physical sciences | Optics | Physics |
4687085 | https://en.wikipedia.org/wiki/Streamflow | Streamflow | Streamflow, or channel runoff, is the flow of water in streams and other channels, and is a major element of the water cycle. It is one runoff component, the movement of water from the land to waterbodies, the other component being surface runoff. Water flowing in channels comes from surface runoff from adjacent hillslopes, from groundwater flow out of the ground, and from water discharged from pipes. The discharge of water flowing in a channel is measured using stream gauges or can be estimated by the Manning equation. The record of flow over time is called a hydrograph. Flooding occurs when the volume of water exceeds the capacity of the channel.
Role in the water cycle
Streams play a critical role in the hydrologic cycle that is essential for all life on Earth. A diversity of biological species, from unicellular organisms to vertebrates, depend on flowing-water systems for their habitat and food resources. Rivers are major aquatic landscapes for all manners of plants and animals. Rivers even help keep the aquifers underground full of water by discharging water downward through their streambeds. In addition to that, the oceans stay full of water because rivers and runoff continually refreshes them. Streamflow is the main mechanism by which water moves from the land to the oceans or to basins of interior drainage.
Sources
Stream discharge is derived from four sources: channel precipitation, overland flow, interflow, and groundwater.
Channel precipitation is the moisture falling directly on the water surface, and in most streams, it adds very little to discharge.
Groundwater enters the streambed where the channel intersects the water table, providing a steady supply of water, termed baseflow, during both dry and rainy periods. Because of the large supply of groundwater available to the streams and the slowness of the response of groundwater to precipitation events, baseflow changes only gradually over time, and it is rarely the main cause of flooding. However, it does contribute to flooding by providing a stage onto which runoff from other sources is superimposed.
Interflow is water that infiltrates the soil and then moves laterally to the stream channel in the zone above the water table. Much of this water is transmitted within the soil, some of it moving within the horizons. Next to baseflow, it is the most important source of discharge for streams in forested lands. Overland flow in heavily forested areas makes negligible contributions to streamflow.
In dry regions, cultivated, and urbanized areas, overland flow or surface runoff is usually a major source of streamflow. Overland flow is a stormwater runoff that begins as thin layer of water that moves very slowly (typically less than 0.25 feet per second) over the ground. Under intensive rainfall and in the absence of barriers such as rough ground, vegetation, and absorbing soil, it can mount up, rapidly reaching stream channels in minutes and causing sudden rises in discharge. The quickest response times between rainfall and streamflow occur in urbanized areas where yard drains, street gutters, and storm sewers collect overland flow and route it to streams straightaway. Runoff velocities in storm sewer pipes can reach 10 to 15 feet per second.
Mechanisms that cause changes in streamflow
Rivers are always moving, which is good for environment, as stagnant water does not stay fresh and inviting very long. There are many factors, both natural and human-induced, that cause rivers to continuously change:
Natural mechanisms
Runoff from rainfall and snowmelt
Evaporation from soil and surface-water bodies
Transpiration by vegetation
Ground-water discharge from aquifers
Ground-water recharge from surface-water bodies
Sedimentation of lakes and wetlands
Formation or dissipation of glaciers, snowfields, and permafrost
Human-induced mechanisms
Surface-water withdrawals and transbasin diversions
River-flow regulation for hydropower and navigation
Construction, removal, and sedimentation of reservoirs and stormwater retention ponds
Stream channelization and levee construction
Drainage or restoration of wetlands
Land use changes such as urbanization that alter rates of erosion, infiltration, overland flow, or evapotranspiration
Wastewater outfalls
Irrigation
Measurement
Streamflow is measured as an amount of water passing through a specific point over time. The units used in the United States are cubic feet per second, while in most other countries cubic meters per second are utilized. There are a variety of ways to measure the discharge of a stream or canal. A stream gauge provides continuous flow over time at one location for water resource and environmental management or other purposes. Streamflow values are better indicators than gage height of conditions along the whole river. Measurements of streamflow are made about every six weeks by United States Geological Survey (USGS) personnel. They wade into the stream to make the measurement or do so from a boat, bridge, or cableway over the stream. For each gaging station, a relation between gage height and streamflow is determined by simultaneous measurements of gage height and streamflow over the natural range of flows (from very low flows to floods). This relation provides the streamflow data from that station. For purposes that do not require a continuous measurement of stream flow over time, current meters or acoustic Doppler velocity profilers can be used. For small streams—a few meters wide or smaller—weirs may be installed.
Approximation
One informal method that provides an approximation of the stream flow termed the orange method or float method is:
Measure a length of stream, and mark the start and finish points. The longest length without changing stream conditions is desired to obtain the most accurate measurement.
Place an orange at the starting point and measure the time for it to reach the finish point with a stopwatch. Repeat this at least three times and average the measurement times.
Express velocity in meters per second. If the measurements were made at midstream (maximum velocity), the mean stream velocity is approximately 0.8 of the measured velocity for rough (rocky) bottom conditions and 0.9 of the measured velocity for smooth (mud, sand, smooth bedrock) bottom conditions.
Monitoring
In the United States, streamflow gauges are funded primarily from state and local government funds. In fiscal year 2008, the USGS provided 35% of the funding for everyday operation and maintenance of gauges. Additionally, USGS uses hydrographs to study streamflow in rivers. A hydrograph is a chart showing, most often, river stage (height of the water above an arbitrary altitude) and streamflow (amount of water, usually in cubic feet per second). Other properties, such as rainfall and water quality parameters can also be plotted.
Forecasting
For most streams especially those with a small watershed, no record of discharge is available. In that case, it is possible to make discharge estimates using the rational method or some modified version of it. However, if chronological records of discharge are available for a stream, a short term forecast of discharge can be made for a given rainstorm using a hydrograph.
Unit hydrograph method
This method involves building a graph in which the discharge generated by a rainstorm of a given size is plotted over time, usually hours or days. It is called the unit hydrograph method because it addresses only the runoff produced by a particular rainstorm in a specified period of time—the time taken for a river to rise, peak, and fall in response to a storm.
Once a rainfall-runoff relationship is established, then subsequent rainfall data can be used to forecast streamflow for selected storms, called standard storms. A standard rainstorm is a high intensity storm of some known magnitude and frequency. One method of unit hydrograph analysis involves expressing the hour by hour or day by day increase in streamflow as a percentage of total runoff. Plotted on a graph, these data from the unit hydrograph for that storm, which represents the runoff added to the pre-storm baseflow.
To forecast the flows in a large drainage basin using the unit hydrograph method would be difficult because in a large basin geographic conditions may vary significantly from one part of the basin to another. This is especially so with the distribution of rainfall because an individual rainstorm rarely covers the basin evenly. As a result, the basin does not respond as a unit to a given storm, making it difficult to construct a reliable hydrograph.
Magnitude and frequency method
For large basins, where unit hydrograph might not be useful and reliable, the magnitude and frequency method is used to calculate the probability of recurrence of large flows based on records of past years' flows. In United States, these records are maintained by the Hydrological Division of the USGS for large streams. For a basin with an area of 5,000 square miles or more, the river system is typically gauged at five to ten places.
The data from each gauging station apply to the part of the basin upstream that location. Given several decades of peak annual discharges for a river, limited projections can be made to estimate the size of some large flow that has not been experienced during the period of record. The technique involves projecting the curve (graph line) formed when peak annual discharges are plotted against their respective recurrence intervals. However, in most cases the curve bends strongly, making it difficult to plot a projection accurately. This problem can be overcome by plotting the discharge and/or recurrence interval data on logarithmic graph paper. Once the plot is straightened, a line can be ruled drawn through the points. A projection can then be made by extending the line beyond the points and then reading the appropriate discharge for the recurrence interval in question.
Relationship to the environment
Runoff of water in channels is responsible for transport of sediment, nutrients, and pollution downstream. Without streamflow, the water in a given watershed would not be able to naturally progress to its final destination in a lake or ocean. This would disrupt the ecosystem. Streamflow is one important route of water from the land to lakes and oceans. The other main routes are surface runoff (the flow of water from the land into nearby watercourses that occurs during precipitation and as a result of irrigation), flow of groundwater into surface waters, and the flow of water from constructed pipes and channels.
Relationship to society
Streamflow confers on society both benefits and hazards. Runoff downstream is a means to collect water for storage in dams for power generation of water abstraction. The flow of water assists transport downstream. A given watercourse has a maximum streamflow rate that can be accommodated by the channel that can be calculated. If the streamflow exceeds this maximum rate, as happens when an excessive amount of water is present in the watercourse, the channel cannot handle all the water, and flooding occurs.
The 1993 Mississippi river flood, the largest ever recorded on the river, was a response to a heavy, long duration spring and summer rainfalls. Early rains saturated the soil over more than a 300,000 square miles of the upper watershed, greatly reducing infiltration and leaving soils with little or no storage capacity. As rains continued, surface depressions, wetlands, ponds, ditches, and farm fields filled with overland flow and rainwater. With no remaining capacity to hold water, additional rainfall was forced from the land into tributary channels and thence to the Mississippi River. For more than a month, the total load of water from hundreds of tributaries exceeded the Mississippi's channel capacity, causing it to spill over its banks onto adjacent floodplains. Where the flood waters were artificially constricted by an engineered channel bordered by constructed levees and unable to spill onto large section of floodplain, the flood levels forced even higher.
| Physical sciences | Hydrology | Earth science |
4688195 | https://en.wikipedia.org/wiki/Astigmatism | Astigmatism | Astigmatism is a type of refractive error due to rotational asymmetry in the eye's refractive power. This results in distorted or blurred vision at any distance. Other symptoms can include eyestrain, headaches, and trouble driving at night. Astigmatism often occurs at birth and can change or develop later in life. If it occurs in early life and is left untreated, it may result in amblyopia.
The cause of astigmatism is unclear, although it is believed to be partly related to genetic factors. The underlying mechanism involves an irregular curvature of the cornea and protective reaction changes in the lens of the eye, called lens astigmatism, that has the same mechanism as spasm of accommodation. Diagnosis is by an eye examination called autorefractor keratometry (objective, allows to see lens and cornea components of astigmatism) and subjective refraction.
Three treatment options are available: glasses, contact lenses, and surgery. Glasses are the simplest. Contact lenses can provide a wider field of vision and fewer artifacts than even double aspheric lenses. Refractive surgery aims to permanently change the shape of the eye and thereby cure astigmatism.
In Europe and Asia, astigmatism affects between 30% and 60% of adults. People of all ages can be affected by astigmatism. Astigmatism was first reported by Thomas Young in 1801.
Signs and symptoms
Although astigmatism may be asymptomatic, higher degrees of astigmatism may cause symptoms such as blurred vision, double vision, squinting, eye strain, fatigue, or headaches. Some research has pointed to the link between astigmatism and higher prevalence of migraine headaches.
Causes
Congenital
The cause of congenital astigmatism is unclear, although it is believed to be partly related to genetic factors. Genetics, based on twin studies, appear to play only a small role in astigmatism as of 2007.
Genome-wide association studies (GWAS) have been used to investigate the genetic foundation of astigmatism. Although no conclusive result has been shown, various candidates have been identified. In a study conducted in 2011 on various Asian populations, variants in the PDGFRA gene on chromosome 4q12 were identified to be associated with corneal astigmatism. A follow-up study in 2013 on the European population, however, found no variant significantly associated with corneal astigmatism at the genome-wide level (single-nucleotide polymorphism rs7677751 at PDGFRA). Facing the inconsistency, a study by Shah and colleagues in 2018 included both populations with Asian and Northern European ancestry. They successfully replicated the previously identified genome-wide significant locus for corneal astigmatism near the PDGFRA gene, with a further success of identifying three novel candidate genes: CLDN7, ACP2, and TNFAIP8L3. Other GWAS studies also provided inconclusive results: Lopes and colleagues identified a susceptibility locus with lead single nucleotide polymorphism rs3771395 on chromosome 2p13.3 in the VAX2 gene (VAX2 plays an important role in the development of the dorsoventral axis of the eye); Li and associates, however, found no consistent or strong genetic signals for refractive astigmatism while suggesting a possibility of widespread genetic co-susceptibility for spherical and astigmatic refractive errors. They also found that the TOX gene region previously identified for spherical equivalent refractive error was the second most strongly associated region. Another recent follow-up study again had identified four novel loci for corneal astigmatism, with two also being novel loci for astigmatism: ZC3H11B (associated with axial length), NPLOC4 (associated with myopia), LINC00340 (associated with spherical equivalent refractive error) and HERC2 (associated with eye color).
Acquired
Astigmatism may also occur following a cataract surgery or a corneal injury. Contraction of the scar due to wound or cataract extraction causes astigmatism due to flattening of the cornea in one direction. In keratoconus, progressive thinning and steepening of the cornea cause irregular astigmatism.
Pathophysiology
Axis of the principal meridian
Regular astigmatism – principal meridians are perpendicular. (The steepest and flattest meridians of the eye are called principal meridians.)
With-the-rule astigmatism – the vertical meridian is steepest (a rugby ball or American football lying on its side).
Against-the-rule astigmatism – the horizontal meridian is steepest (a rugby ball or American football standing on its end).
Oblique astigmatism – the steepest curve lies in between 120 and 150 degrees and 30 and 60 degrees.
Irregular astigmatism – principal meridians are not perpendicular.
In with-the-rule astigmatism, the eye has too much "plus" cylinder in the horizontal axis relative to the vertical axis (i.e., the eye is too "steep" along the vertical meridian relative to the horizontal meridian). Vertical beams of light focus in front (anterior) to horizontal beams of light, in the eye. This problem may be corrected using spectacles which have a "minus" cylinder placed on this horizontal axis. The effect of this will be that when a vertical beam of light in the distance travels towards the eye, the "minus" cylinder (which is placed with its axis lying horizontally – meaning in line with the patient's horizontal meridian relative to the excessively steep vertical meridian) will cause this vertical beam of light to slightly "diverge", or "spread out vertically", before it reaches the eye. This compensates for the fact that the patient's eye converges light more powerfully in the vertical meridian than the horizontal meridian. Hopefully, after this, the eye will focus all light on the same location at the retina, and the patient's vision will be less blurred.
In against-the-rule astigmatism, a plus cylinder is added in the horizontal axis (or a minus cylinder in the vertical axis).
Axis is always recorded as an angle in degrees, between 0 and 180 degrees in a counter-clockwise direction. Both 0 and 180 degrees lie on a horizontal line at the level of the center of the pupil, and as seen by an observer, 0 lies on the right of both the eyes.
Irregular astigmatism, which is often associated with prior ocular surgery or trauma, is also a common naturally occurring condition. The two steep hemimeridians of the cornea, 180° apart in regular astigmatism, may be separated by less than 180° in irregular astigmatism (called nonorthogonal irregular astigmatism); and/or the two steep hemimeridians may be asymmetrically steep—that is, one may be significantly steeper than the other (called asymmetric irregular astigmatism). Irregular astigmatism is quantified by a vector calculation called topographic disparity.
Focus of the principal meridian
With accommodation relaxed:
Simple astigmatism
Simple hyperopic astigmatism – first focal line is on the retina, while the second is located behind the retina.
Simple myopic astigmatism – first focal line is in front of the retina, while the second is on the retina.
Compound astigmatism
Compound hyperopic astigmatism – both focal lines are located behind the retina.
Compound myopic astigmatism – both focal lines are located in front of the retina.
Mixed astigmatism – focal lines are on both sides of the retina (straddling the retina).
Throughout the eye
Astigmatism, whether it is regular or irregular, is caused by some combination of external (corneal surface) and internal (posterior corneal surface, human lens, fluids, retina, and eye-brain interface) optical properties. In some people, the external optics may have the greater influence, and in other people, the internal optics may predominate. Importantly, the axes and magnitudes of external and internal astigmatism do not necessarily coincide, but it is the combination of the two that by definition determines the overall optics of the eye. The overall optics of the eye are typically expressed by a person's refraction; the contribution of the external (anterior corneal) astigmatism is measured through the use of techniques such as keratometry and corneal topography. One method analyzes vectors for planning refractive surgery such that the surgery is apportioned optimally between both the refractive and topographic components.
Diagnosis
A number of tests are used during eye examinations to determine the presence of astigmatism and to quantify its amount and axis. A Snellen chart or other eye charts may initially reveal reduced visual acuity. A keratometer may be used to measure the curvature of the steepest and flattest meridians in the cornea's front surface. Corneal topography may also be used to obtain a more accurate representation of the cornea's shape. An autorefractor or retinoscopy may provide an objective estimate of the eye's refractive error and the use of Jackson cross cylinders in a phoropter or trial frame may be used to subjectively refine those measurements. An alternative technique with the phoropter requires the use of a "clock dial" or "sunburst" chart to determine the astigmatic axis and power. A keratometer may also be used to estimate astigmatism by finding the difference in power between the two primary meridians of the cornea. Javal's rule can then be used to compute the estimate of astigmatism.
A method of astigmatism analysis by Alpins may be used to determine both how much surgical change of the cornea is needed and after surgery to determine how close treatment was to the goal.
Another rarely used refraction technique involves the use of a stenopaeic slit (a thin slit aperture) where the refraction is determined in specific meridians – this technique is particularly useful in cases where the patient has a high degree of astigmatism or in refracting patients with irregular astigmatism.
Classification
There are three primary types of astigmatism: myopic astigmatism, hyperopic astigmatism, and mixed astigmatism. Cases can be classified further, such as regular or irregular and lenticular or corneal.
Treatment
Astigmatism may be corrected with eyeglasses, contact lenses, or refractive surgery. Glasses are the simplest and safest, although contact lenses can provide a wider field of vision. Refractive surgery can eliminate the need to wear corrective lenses altogether by permanently changing the shape of the eye but, like all elective surgery, comes with both greater risk and expense than the non-invasive options. Various considerations involving eye health, refractive status, and lifestyle determine whether one option may be better than another. In those with keratoconus, certain contact lenses often enable patients to achieve better visual acuity than eyeglasses. Once only available in a rigid, gas-permeable form, toric lenses are now also available as soft lenses.
In older people, astigmatism can also be corrected during cataract surgery. This can either be done by inserting a toric intraocular lens or by performing special incisions (limbal relaxing incisions). Toric intraocular lenses probably provide a better outcome with respect to astigmatism in these cases than limbal relaxing incisions.
Toric intraocular lenses can additionally be used in patients with complex ophthalmic history, such as previous ophthalmic surgery. In such complex cases, toric intraocular lenses seem to be as effective as in non-complex cases for correction of concurrent corneal astigmatism.
Epidemiology
In 2019, the World Health Organization reported that 123.7 million people worldwide were affected by uncorrected refracting errors, including astigmatism. A compilation of many systematic reviews found that there was an 8-62% prevalence of astigmatism among adults, with an estimated prevalence of 40% worldwide. The country with the highest reported prevalence among the compilation of systematic reviews is China at 62%. The prevalence of astigmatism increases with age due to changes in refractive index gradients. According to an American study, nearly three in ten children (28.4%) between the ages of five and seventeen have astigmatism. A Brazilian study published in 2005 found that 34% of the students in one city were astigmatic.
Studies have shown that infants in their first few months have a high prevalence of astigmatism due to a steep cornea. The steepest corneas are found in infants with low birth weights and post-conceptional age. By the age of four, the prevalence of astigmatism has reduced as the cornea flattens. The cornea remains mostly stable during adulthood, and then steepens again in older adulthood (40+ years).
Mild astigmatism has a higher prevalence than moderate and significant astigmatisms and increased until the age of 70, while moderate and significant astigmatisms showed an increase in prevalence after the age of 70. Of the levels of astigmatism, mild astigmatism is most prevalent, making up about 82% of the total reported astigmatisms.
With-the-rule astigmatism (from studies with differing age groups) has a prevalence range of 4 to 98% globally. The prevalence range for against-the-rule astigmatism (from studies with differing age groups) is from 1 to 58%. For oblique astigmatism, the prevalence range is from 2 to 61%. With-the-rule astigmatism is more prevalent in young adults, and over time, the prevalence shifts to be mostly against-the-rule astigmatism. A Polish study published in 2005 revealed "with-the-rule astigmatism" may lead to the onset of myopia.
The main cause of astigmatism is changes in the curvature of the cornea. When left untreated, astigmatism causes people to have a lower vision-related quality of life. Some factors that lead to this are a decrease in vision quality and an increase in glare and haloes. People with astigmatism have more difficulty with night driving and can have a decreased productivity due to errors. However, there are many ways to help correct astigmatisms: The use of glasses or contacts, Toric intraocular lenses, Toric implantable Collamer lenses, and/or corneal refractive surgery have been shown to correct astigmatisms.
History
As a student, Thomas Young discovered that he had problems with one eye in 1793. In the following years, he did research on his vision problems. He presented his findings in a Bakerian Lecture in 1801.
Independent from Young, George Biddell Airy discovered the phenomenon of astigmatism on his own eye. Airy presented his observations on his own eye in February 1825 at the Cambridge Philosophical Society. Airy produced lenses to correct his vision problems by 1825, while other sources put this into 1827 when Airy obtained cylindrical lenses from an optician from Ipswich. The name for the condition was given by William Whewell.
By the 1860s, astigmatism was a well established concept in ophthalmology, and chapters in books described the discovery of astigmatism.
In 1849, Irish English physicist and mathematician George Stokes invented Stokes lens to detect astigmatism. In 1887, American ophthalmologist Edward Jackson revised the Stokes lens concept and made a cross cylinder lens to refine power and axis of astigmatism. In 1907, Jackson described determination of the axis of a correcting cylinder in astigmatism using a cross cylinder.
| Biology and health sciences | Disabilities | Health |
4689328 | https://en.wikipedia.org/wiki/TNT%20equivalent | TNT equivalent | TNT equivalent is a convention for expressing energy, typically used to describe the energy released in an explosion. The is a unit of energy defined by convention to be (), which is the approximate energy released in the detonation of a metric ton (1,000 kilograms) of TNT. In other words, for each gram of TNT exploded, (or 4184 joules) of energy are released.
This convention intends to compare the destructiveness of an event with that of conventional explosive materials, of which TNT is a typical example, although other conventional explosives such as dynamite contain more energy.
Kiloton and megaton
The "kiloton (of TNT equivalent)" is a unit of energy equal to 4.184 terajoules (). A kiloton of TNT can be visualized as a cube of TNT on a side.
The "megaton (of TNT equivalent)" is a unit of energy equal to 4.184 petajoules ().
The kiloton and megaton of TNT equivalent have traditionally been used to describe the energy output, and hence the destructive power, of a nuclear weapon. The TNT equivalent appears in various nuclear weapon control treaties, and has been used to characterize the energy released in asteroid impacts.
Historical derivation of the value
Alternative values for TNT equivalency can be calculated according to which property is being compared and when in the two detonation processes the values are measured.
Where for example the comparison is by energy yield, an explosive's energy is normally expressed for chemical purposes as the thermodynamic work produced by its detonation. For TNT this has been accurately measured as 4,686 J/g from a large sample of air blast experiments, and theoretically calculated to be 4,853 J/g.
However even on this basis, comparing the actual energy yields of a large nuclear device and an explosion of TNT can be slightly inaccurate. Small TNT explosions, especially in the open, don't tend to burn the carbon-particle and hydrocarbon products of the explosion. Gas-expansion and pressure-change effects tend to "freeze" the burn rapidly. A large open explosion of TNT may maintain fireball temperatures high enough so that some of those products do burn up with atmospheric oxygen.<ref
name="Needham">
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Such differences can be substantial. For safety purposes a range as wide as has been stated for a gram of TNT upon explosion. Thus one can state that a nuclear bomb has a yield of 15 kt (), but the explosion of an actual pile of TNT may yield (for example) due to additional carbon/hydrocarbon oxidation not present with small open-air charges.
These complications have been sidestepped by convention. The energy released by one gram of TNT was arbitrarily defined as a matter of convention to be 4,184 J, which is exactly one kilocalorie.
Conversion to other units
1 ton of TNT equivalent is approximately:
calories
joules
British thermal units
foot-pounds
kilowatt-hours
electronvolts
Examples
Relative effectiveness factor
The relative effectiveness factor (RE factor) relates an explosive's demolition power to that of TNT, in units of the TNT equivalent/kg (TNTe/kg). The RE factor is the relative mass of TNT to which an explosive is equivalent: The greater the RE, the more powerful the explosive.
This enables engineers to determine the proper masses of different explosives when applying blasting formulas developed specifically for TNT. For example, if a timber-cutting formula calls for a charge of 1 kg of TNT, then based on octanitrocubane's RE factor of 2.38, it would take only 1.0/2.38 (or 0.42) kg of it to do the same job. Using PETN, engineers would need 1.0/1.66 (or 0.60) kg to obtain the same effects as 1 kg of TNT. With ANFO or ammonium nitrate, they would require 1.0/0.74 (or 1.35) kg or 1.0/0.32 (or 3.125) kg, respectively.
Calculating a single RE factor for an explosive is, however, impossible. It depends on the specific case or use. Given a pair of explosives, one can produce 2× the shockwave output (this depends on the distance of measuring instruments) but the difference in direct metal cutting ability may be 4× higher for one type of metal and 7× higher for another type of metal. The relative differences between two explosives with shaped charges will be even greater. The table below should be taken as an example and not as a precise source of data.
Nuclear examples
| Physical sciences | Energy | Basics and measurement |
4691461 | https://en.wikipedia.org/wiki/Rotterdam%20Metro | Rotterdam Metro | The Rotterdam Metro () is a rapid transit system operated in Rotterdam, Netherlands and surrounding municipalities by RET. The first line, called Noord – Zuidlijn (North – South line) opened in 1968 and ran from Centraal Station to Zuidplein, crossing the river Nieuwe Maas in a tunnel. It was the first metro system to open in the Netherlands. At the time it was also one of the shortest metro lines in the world with a length of only .
In 1982 a second line was opened, the Oost – Westlijn (East – West line), running between Capelsebrug and Coolhaven stations. In the late 1990s, the lines were named after two historic Rotterdam citizens, the Erasmus Line (North – South) after Desiderius Erasmus and the Caland Line (East – West) after Pieter Caland. As of December 2009, these names were dropped again in favour of a combination of letters and colours, to emphasise and clarify the difference between the separate branches, especially of the former East – West line.
Lines
Lines A and B
In the northeast of Rotterdam, Lines A and B branch to Binnenhof (Line A) and to Nesselande (Line B). The latter has been extended since September 2005; before that date, this line terminated at De Tochten.
North of Capelsebrug station and west of Schiedam Centrum station, with the exception of the De Tochten-Nesselande section, Lines A and B have some level crossings (with priority), and could therefore be called light rail instead of metro. These sections also have overhead wires, while most of the system has a third rail (the other exception is Line E (RandstadRail) to The Hague). However, the term light rail is not used in Rotterdam; most people just call these branches metro.
As of 30 September 2019 Line B is connected to the Schiedam–Hoek van Holland railway line, extending the metro network to Hook of Holland (Hoek van Holland), while line A operates on this line as far as Vlaardingen West, starting 1 November 2019.
Line C
At Capelsebrug, Line C branches off the main East-West section to De Terp in Capelle aan den IJssel. Until November 2002, the Calandlijn (now lines A, B and C) terminated in the west of Rotterdam, at Marconiplein. On 4 November 2002 an extension through the city of Schiedam towards Spijkenisse was opened. The extension included four new stations in Schiedam (including Schiedam Centrum station) and one in Pernis. Line C joins Line D at the Tussenwater station in Hoogvliet. Line A and B branch off to the Schiedam-Hoek van Holland railway, while Line C trains continue and, like those on the Line D, terminate at De Akkers station in Spijkenisse.
Line D
Line D runs from Rotterdam Centraal via Beurs, Slinge, Rhoon, Tussenwater, and Spijkenisse Centrum towards De Akkers.
Line D intersects with Lines A, B and C at Beurs station. Before the connection with Line E at Rotterdam Centraal was realized in December 2011, some Line D trains terminated at Slinge during rush hours.
Line E
Line E is a direct conversion of the former Hofpleinlijn from a railway line to a rapid transit line in 2006. The section of the route between Laan van NOI and Leidschendam-Voorburg is shared with light rail vehicles on two routes from the Hague tram network heading towards Zoetermeer via the Zoetermeer Stadslijn. All of these services fall under the RandstadRail branding.
At the time of opening, the old Hofplein station was temporarily kept as the line's southern terminus. On 17 August 2010, however, a new tunnel opened, which connected the metro station at Rotterdam Centraal via a new tunnel and new Blijdorp station with the existing tracks near Melanchthonweg station.
For the next year, work was in progress to connect Line D to Line E at Rotterdam Centraal station. Since the completion of this project in December 2011, all trains coming from Den Haag Centraal terminate at Slinge (these are Line E trains), while line D continues in service between De Akkers and Rotterdam Centraal.
The line's northern terminus used the former railway platforms at Den Haag Centraal, which had been inherited from NS and in use since 1975, having been moved from the previous terminus at Den Haag Hollands Spoor. These platforms remained in use until 12 February 2016. The line was then closed between here and Laan Van NOI while construction was completed of a new station built on an elevated viaduct adjacent to the railway platforms. This opened on 22 August 2016.
Network map
Overview
In 2015, there were about 175,000 daily riders of Lines A, B and C. There were 145,000 on Lines D and E.
Future extensions
RET plans to build a connecting line from Kralingse Zoom station to the new Feijenoord City housing area, and then on to Zuidplein, Charlois and Rotterdam Central station; along with plans to convert two tracks of four heavy-rail tracks between Rotterdam Central and Dordrecht to metro operation. By adding extra stations and operating trains at two-minute interval RET expects to achieve a significant increase in traffic. Another ambition is to automate the metro to achieve 90-second headways.
Rolling stock
Trains series 5300, 5400, 5500, 5600 and 5700 are Bombardier Flexity Swifts. The series 5500 trains, made between 2007 and 2009, were built for the new RandstadRail Line E. The 5601-5642 trains were built to replace older Duewag stock (series 5200). In 2013, RET announced that it ordered 16 additional vehicles of SG3 stock to run on the Hoekse Line extension. A further six vehicles were ordered to increase capacity on the Randstadrail branch in 2015. Delivery of these 22 vehicles, called HSG3, took place between 2015 and 2017.
Traction power
Trains run on 750 volts DC power which is supplied through a bottom-contact third rail throughout most of the system. There are multiple spring-loaded contact shoes on both sides of the vehicle, which are loaded and unloaded automatically due to the slanted edges of third rail ends. This allows the rail to be installed on either side of the track, a necessity around points and station platforms. There is sufficient overlap between the two rails on either end to avoid a "gapped" train, a situation where none of the shoes are in contact with the live rail. To reduce the risk of electrocution, the rail consists of a sturdy yellow insulating material, with the live current carried on a thick metal strip on the bottom side. This also guards against grime (such as from fallen autumn leaves) reducing or preventing electrical contact.
Three lines do however have sections that use overhead wires. These are lines A (towards Binnenhof), B (towards Nesselande), and E (towards Den Haag Centraal). On lines A and B, trains raise or lower their pantographs while the vehicle is in motion just east of Capelsebrug station, while on line E this happens while stationary at Melanchtonweg station (this leads to the only level crossing with third rail in the country being at the Kleiweg just outside the tunnel heading to Blijdorp station). Note that Line B trains switch back to third rail for the final leg of the journey, from the penultimate station De Tochten to Nesselande. The western extension of lines A and B to Hook of Holland also use overhead power as they are converted directly from the existing railway line.
The sections of the metro that use overhead wires are called sneltram (light rail) by locals, as they include several protected level crossings at street level, which trains pass through with priority, as in a conventional railway line. For this reason, trains with a pantograph (series 5200 and 5400) are equipped with turn signals just like any road vehicle. This makes it easy to see the difference between series 5300 and series 5400 Bombardier-built trains. Series 5500 and 5600 trains are also equipped, although the former is normally used on line E only as they carry the RandstadRail branding and livery while the latter carries that of R-net.
| Technology | Netherlands | null |
4692188 | https://en.wikipedia.org/wiki/Two-toed%20sloth | Two-toed sloth | Choloepus is a genus of xenarthran mammals from Central and South America within the monotypic family Choloepodidae, consisting of two-toed sloths, sometimes also called two-fingered sloths. The two species of Choloepus (which means "lame foot" in Ancient Greek), Linnaeus's two-toed sloth (Choloepus didactylus) and Hoffmann's two-toed sloth (Choloepus hoffmanni), were formerly believed on the basis of morphological studies to be the only surviving members of the sloth family Megalonychidae, but have now been shown by molecular results to be closest to extinct ground sloths of the family Mylodontidae.
Extant species
Evolution
A study of retrovirus and mitochondrial DNA suggests that C. didactylus and C. hoffmani diverged 6 to 7 million years ago. Furthermore, based on cytochrome c oxidase subunit I sequences, a similar divergence date ( years ago) between the two populations of C. hofmanni separated by the Andes has been reported. Their ancestors evolved with marine vertebrae, the three toed-sloth and the manatee are the only other mammals with similar vertebrae.
Relation to the three-toed sloth
Both types of sloth tend to occupy the same forests; in most areas, a particular species of the somewhat smaller and generally slower-moving three-toed sloth (Bradypus) and a single species of the two-toed type will jointly predominate. Although similar in overall appearance, the relationship between the two genera is not close. Recent phylogenetic analyses support analysis of morphological data from the 1970s and 1980s, indicating the two genera are not closely related and adapted to their arboreal lifestyles independently. It was unclear from this work from which ground-dwelling sloth taxa the three-toed sloths evolved. Based on the morphological comparisons, it was thought the two-toed sloths nested phylogenetically within one of the divisions of Caribbean sloths. Though data has been collected on over 33 different species of sloths by analyzing bone structures, many of the relationships between clades on a phylogenetic tree were unclear.
Much of the morphological evidence to support the hypothesis of diphyly has been based on the structure of the inner ear. Most morphological studies have concluded that convergent evolution is the mechanism that resulted in today's two genera of tree sloths. This means that the extant genera evolved analogous traits, such as locomotion methods, size, habitat, and many other traits independently from one another as opposed to from their last common ancestor. This makes tree sloths “one of the most striking examples of convergent evolution known among mammals”.
Recently obtained molecular data from collagen and mitochondrial DNA sequences fall in line with the diphyly (convergent evolution) hypothesis, but have overturned some of the other conclusions obtained from morphology. These investigations consistently place two-toed sloths close to mylodontids and three-toed sloths within Megatherioidea, close to Megalonyx, megatheriids and nothrotheriids. They make the previously recognized family Megalonychidae polyphyletic, with both two-toed sloths and the Caribbean sloths being moved out of that family and away from Megalonyx. Caribbean sloths are placed in a separate, basal branch of the sloth evolutionary tree.
Characteristics
The name "two-toed sloth" was intended to describe an anatomical difference between the genera Choloepus and Bradypus, but does so in a potentially misleading way. Members of Choloepus have two digits on their forelimbs (the thoracic limbs) and three digits on their hindlimbs (the pelvic limbs), while members of Bradypus have three digits on all limbs. Although the term "two-fingered" sloth is arguably less misleading, the shorter "two-toed" is much more widely used.
Members of Choloepus are larger than three-toed sloths, having a body length of , and weighing . Other distinguishing features include a more prominent snout, longer fur, and the absence of a tail.
Behaviour
Two-toed sloths spend most of their lives hanging upside down from trees. They cannot walk, so they pull hand-over-hand to move around, which is at an extremely slow rate. Almost all of their movement comes from this suspended upside down position, at a higher degree than even three-toed sloths. As a result, they tend to gravitate towards less vertical portions of trees. Being predominantly nocturnal, their fur, which grows greenish algae to blend in, is their main source of protection. Their body temperatures depend at least partially on the ambient temperature; they cannot shiver to keep warm, as other mammals do, because of their unusually low metabolic rates and reduced musculature. Two-toed sloths also differ from three-toed sloths in their climbing behavior, preferring to descend head first.
Lifecycle
Two-toed sloths have a gestation period of six months to a year, depending on the species. Their ovarian cycle lasts around 31 to 33 days, independently of the seasons but dependent on the species. The mother gives birth to a single young, while hanging upside down. The young are born with claws, and are weaned after about a month, although they will remain with the mother for several more months, and do not reach maturity until the age of three years, in the case of females, or four to five years, in the case of males. During natal dispersion, two-toed sloths prefer tropical forests over other types of habitat, often using riparian forest buffers to disperse. Although they also occupy shade-grown cacao plantations, they avoid open pastures.
Feeding
They eat primarily leaves, but also shoots, fruits, nuts, berries, bark, some native flowers, and even some small vertebrates. In addition, when they cannot find food, they have been known to eat the algae that grow on their fur for nutrients. They have also been observed using mineral licks. They have large, four-chambered stomachs, which help to ferment the large amount of plant matter they eat. Food can take up to a month to digest due to their slow metabolism. Depending on when in the excretion cycle a sloth is weighed, urine and feces may account for up to 30% of the animal's body weight. They get most of their fluids from water in the leaves that they eat but sloths have also been observed drinking directly from rivers.
Dentition and skeleton
Two-toed sloths have a reduced, ever growing dentition, with no incisors or true canines, which overall lacks homology with the dental formula of other mammals. Their first tooth is very canine-like in shape and is referred to as a caniniform. It is used for tearing small chunks off of their food, as well as for defense against predators. It is separated from the other teeth, or molariforms, by a diastema. The molariforms are used specifically for grinding and are mortar and pestle-like in appearance and function. Thus, they can grind food for easier digestibility, which takes the majority of their energy. The dental formula of two-toed sloths is: (unau)
Two-toed sloths are unusual among mammals in possessing as few as five cervical vertebrae, which may be due to mutations in the homeotic genes. All other mammals have seven cervical vertebrae, other than the three-toed sloth and the manatee.
Musculature
Two-toed sloths generally have similar musculature to that of other mammals. This includes their zygomaticus muscles, their superficial masseter, their deep masseter, and their medial and lateral pterygoids. Additionally, a specific section of their anterior temporalis is arranged vertically, to allow them to sharpen their caniniform teeth. They tend to have stronger flexor muscles in their fore- and hindlimbs, as well as their shoulders.
| Biology and health sciences | Xenarthra | Animals |
4692632 | https://en.wikipedia.org/wiki/Portrait%20photography | Portrait photography | Portrait photography, or portraiture, is a type of photography aimed toward capturing the personality of a person or group of people by using effective lighting, backdrops, and poses. A portrait photograph may be artistic or clinical. Frequently, portraits are commissioned for special occasions, such as weddings, school events, or commercial purposes. Portraits can serve many purposes, ranging from usage on a personal web site to display in the lobby of a business.
History
The relatively low cost of the daguerreotype in the middle of the 19th century and the reduced sitting time for the subject, though still much longer than now, led to a general rise in the popularity of portrait photography over painted portraiture. The style of these early works reflected the technical challenges associated with long exposure times and the painterly aesthetic of the time. Hidden mother photography, in which portrait photographs featured young children's mothers hidden in the frame to calm them and keep them still, arose from this difficulty. Subjects were generally seated against plain backgrounds, lit with the soft light of an overhead window, and whatever else could be reflected with mirrors.
Advances in photographic technology since the daguerreotype spawned more advanced techniques, allowed photographers to capture images with shorter exposure times, and work outside a studio environment.
Lighting for portraiture
There are many techniques available to light a subject's face.
Three-point lighting
Three-point lighting is one of the most common lighting setups. It is traditionally used in a studio, but photographers may use it on-location in combination with ambient light. This setup uses three lights, the key light, fill light, and back light, to fully bring out details and the three-dimensionality of the subject's features.
Key light
The key light, also known as the main light, is placed either to the left, right, or above the subject's face, typically 30 to 60 degrees from the camera. The purpose of the key light is to give shape to and emphasize particular features of the subject. The distance of the key light from the camera controls the falloff of the light and profoundness of shadows.
Fill light
The fill light, also known as the secondary main light, is typically placed opposite the key light. For example, if the key light is placed 30 degrees camera-left, the fill light will be placed 30 degrees camera-right. The purpose of a fill light is to combat strong shadows created by the main light. Intensity of the fill light may be equal to the main light to eliminate shadows completely, or less intense to simply lessen shadows. Sometimes, the purpose of a fill light may be served by a reflector rather than an actual light.
Back light
The back light, also known as a hair light, helps separate a subject from its background and emphasize hair. In some cases, photographers may use a hair light to create lens flare or other artistic effects.
High-key and low-key lighting
High-key
High-key lighting is a technique used to result in an image that is mostly free of shadows and has a background brighter than the subject. High-key lighting typically involves use of all three lights (or more) in the three-point lighting setup.
Low-key
Low-key lighting is a technique used to result in an image where only part of the subject is lit, has dark shadows, and a background darker than the subject. Low-key lighting typically involves use of just one light in the three-point lighting setup (although sometimes two).
Butterfly lighting
Butterfly lighting uses only two lights. The key light is placed directly in front of the subject above the camera (or slightly to one side), and a bit higher than the key light in a three-point lighting setup. The second light (more often a reflector rather than an actual light) is placed as a fill directly below the camera (or slightly to the opposite side).
This lighting may be recognized by the strong light falling on the forehead, the bridge of the nose, the upper cheeks, and by the distinct shadow below the nose that often looks rather like a butterfly and thus, provides the name for this lighting technique.
Butterfly lighting was a favourite of famed Hollywood portraitist George Hurrell, which is why this style of lighting is often called Paramount lighting.
Other lighting equipment
Most lights used in modern photography are a flash of some sort. The lighting for portraiture is typically diffused by bouncing it from the inside of an umbrella, or by using a soft box. A soft box is a fabric box, encasing a photo strobe head, one side of which is made of translucent fabric. This provides a softer lighting for portrait work and is often considered more appealing than the harsh light often cast by open strobes. Hair and background lights are usually not diffused. It is more important to control light spillage to other areas of the subject. Snoots, barn doors and flags or gobos help focus the lights exactly where the photographer wants them. Background lights are sometimes used with color gels placed in front of the light to create colored backgrounds.
Devices, tools, or accessories employed in photography, videography, and cinematography to shape, control, alter, direct, block, blackout, or otherwise affect light emitted from a light source, which may be natural or artificial light are called Light Modifiers. By altering the quality, direction, intensity, color, and or other attributes and characteristics of light, light modifiers enabling photographers to achieve specific effects or moods in their images, as well as shoot at locations and times that would not be possible without ability to modify light. Light modifiers come in various categories and types, each with its own unique characteristics and applications. They can be freestanding, placed on stands, handheld, hung, fit over a camera lens, etc. Modifiers can be collapsible and portable and/or rigid and stationary.
Windowlight portraiture
Windows as a source of light for portraits have been used for decades before artificial sources of light were discovered. According to Arthur Hammond, amateur and professional photographers need only two things to light a portrait: a window and a reflector. Although window light limits options in portrait photography compared to artificial lights it gives ample room for experimentation for amateur photographers. A white reflector placed to reflect light into the darker side of the subject's face, will even the contrast. Shutter speeds may be slower than normal, requiring the use of a tripod, but the lighting will be beautifully soft and rich.
The best time to take window light portrait is considered to be early hours of the day and late hours of afternoon when light is more intense on the window. Curtains, reflectors, and intensity reducing shields are used to give soft light. While mirrors and glasses can be used for high key lighting. At times colored glasses, filters and reflecting objects can be used to give the portrait desired color effects. The composition of shadows and soft light gives window light portraits a distinct effect different from portraits made from artificial lights.
While using window light, the positioning of the camera can be changed to give the desired effects. Such as positioning the camera behind the subject can produce a silhouette of the individual while being adjacent to the subject give a combination of shadows and soft light. And facing the subject from the same point of light source will produce high key effects with least shadows.
Styles of portraiture
There are many different techniques for portrait photography. Often it is desirable to capture the subject's eyes and face in sharp focus while allowing other less important elements to be rendered in a soft focus. At other times, portraits of individual features might be the focus of a composition such as the hands, eyes or part of the subject's torso.
Head shots have become a popular style within portrait photography, particularly in the entertainment industry, where they are commonly used to showcase a subject's facial features and expressions.
Approaches to portraiture
There are essentially four approaches that can be taken in photographic portraiture—the constructionist, environmental, candid, and creative approach. Each has been used over time for different reasons be they technical, artistic or cultural.
The constructionist approach is when the photographer constructs an idea around the subject. It is the approach used in most studio and social photography. It is also used extensively in advertising and marketing when an idea has to be put across.
The environmental approach depicts the subject in their environment. They are often shown as doing something which relates directly to the subject.
The candid approach is where people are photographed without their knowledge going about their daily business. Whilst this approach taken by the paparazzi has been criticized, less invasive and exploitative candid photography has given the world important images of people in various situations and places over the last century. The images of Parisians by Doisneau and Cartier-Bresson demonstrate this approach. As with environmental photography, candid photography is important as a historical source of information about people.
The creative approach is where manipulation of the image is used to change the final output.
Lenses
Lenses used in portrait photography are classically fast, medium telephoto lenses, though any lens may be used, depending on artistic purposes. See Canon EF Portrait Lenses for Canon lenses in this style; other manufacturers feature similar ranges. The first dedicated portrait lens was the Petzval lens developed in 1840 by Joseph Petzval. It had a relatively narrow field of view of 30 degrees, a focal length of 150 mm, and a fast f-number in the ƒ/3.3-3.7 range.
Classic focal length is in the range 80–135 mm on 135 film format and about 150-400mm on large format, which historically is first in photography. Such a field of view provides a flattening perspective distortion when the subject is framed to include their head and shoulders. Wider angle lenses (shorter focal length) require that the portrait be taken from closer (for an equivalent field size), and the resulting perspective distortion yields a relatively larger nose and smaller ears, which is considered unflattering and imp-like. Wide-angle lenses – or even fisheye lenses – may be used for artistic effect, especially to produce a grotesque image. Conversely, longer focal lengths yield greater flattening because they are used from further away. This makes communication difficult and reduces rapport. They may be used, however, particularly in fashion photography, but longer lengths require a loudspeaker or walkie-talkie to communicate with the subject or assistants. In this range, the difference in perspective distortion between 85mm and 135mm is rather subtle; see for examples and analysis.
Speed-wise, fast lenses (wide aperture) are preferred, as these allow shallow depth of field (blurring the background), which helps isolate the subject from the background and focus attention on them. This is particularly useful in the field, where one does not have a back drop behind the subject, and the background may be distracting. The details of bokeh in the resulting blur are accordingly also a consideration; some lenses, in particular the "DC" (Defocus Control) types by Nikon, are designed to give the photographer control over this aspect, by providing an additional ring acting only on the quality of the bokeh, without influencing the foreground (hence, these are not soft-focus lenses). However, extremely wide apertures are less frequently used, because they have a very shallow depth of field and thus the subject's face will not be completely in focus.
Conversely, in environmental portraits, where the subject is shown in their environment, rather than isolated from it, background blur is less desirable and may be undesirable, and wider angle lenses may be used to show more context.
Finally, soft focus (spherical aberration) is sometimes a desired effect, particularly in glamour photography where the "gauzy" look may be considered flattering. The Canon EF 135mm 2.8 with Softfocus is an example of a lens designed with a controllable amount of soft focus.
Most often a prime lens will be used, both because the zoom is not necessary for posed shots (and primes are lighter, cheaper, faster, and higher quality), and because zoom lenses can introduce highly unflattering geometric distortion (barrel distortion or pincushion distortion). However, zoom lenses may be used, particularly in candid shots or to encourage creative framing.
Portrait lenses are often relatively inexpensive, because they can be built simply, and are close to the normal range. The cheapest portrait lenses are normal lenses (50 mm), used on a cropped sensor. For example, the Canon EF 50mm f/1.8 II is the least expensive Canon lens, but when used on a 1.6× cropped sensor yields an 80mm equivalent focal length, which is at the wide end of portrait lenses.
Mobile portraiture
The documentary I Am Chicago was an experiment in mobile full-body portraiture, using natural light and a moving truck as a studio. The project aimed to break down traditional barriers of access to the art form.
Senior portraits
In North America, senior portraits are formal portraits taken of students near the end of their senior year of high school. Senior portraits are often included in graduation announcements or are given to friends and family. They are also used in yearbooks and are usually rendered larger than their underclassmen counterparts and are often featured in color, even if the rest of the yearbook is mostly reproduced in black and white. In some schools the requirements are strict regarding the choice of photographer or in the style of portraiture, with only traditional-style portraits being acceptable. Many schools choose to contract one photographer for their yearbook portraits, while other schools allow many different photographers to submit yearbook portraits. Senior portraits have become a cultural rite of passage in the United States, representing a momentous achievement in a young person's life and serving as a tangible reminder of their high school years for years to come.
Traditional
Formal senior portraits date back at least to the 1880s in America. Some traditional senior portrait sittings include a cap and gown and other changes of clothing, portrait styles and poses.
In some schools a portrait studio is invited to the school to ensure all senior portraits (for the yearbook) are similar in pose and style, and so that students who cannot afford to purchase these portraits on their own or choose not to purchase portraits will appear in the yearbook the same as other students. Other schools allow students to choose a studio and submit portraits on their own.
Modern
Modern senior portraits may include virtually any pose or clothing choice within the limits of good taste. Students often appear with pets, student athletes pose in letterman jackets or their playing uniforms, and many choose fashion photography. Outdoor photos are popular at locations that are scenic or important to the senior. Picture proofs are usually available to view online the next day which are lower quality, unedited and often with a watermark of the studio.
| Technology | Photography | null |
4693311 | https://en.wikipedia.org/wiki/Thiosulfate | Thiosulfate | Thiosulfate (IUPAC-recommended spelling; sometimes thiosulphate in British English) is an oxyanion of sulfur with the chemical formula . Thiosulfate also refers to the compounds containing this anion, which are the salts of thiosulfuric acid, such as sodium thiosulfate and ammonium thiosulfate . Thiosulfate salts occur naturally. Thiosulfate rapidly dechlorinates water, and is used to halt bleaching in the paper-making industry. Thiosulfate salts are mainly used for dyeing in textiles, and bleaching of natural substances.
Structure and bonding
The thiosulfate ion is tetrahedral at the central S atom. The thiosulfate ion has C3v symmetry. The external sulfur atom has a valence of 2 while the central sulfur atom has a valence of 6. The oxygen atoms have a valence of 2. The S-S distance of about 201 pm in sodium thiosulphate is appropriate for a single bond. The S-O distances are slightly shorter than the S-O distances in sulfate.
For many years, the oxidation states of the sulfur atoms in the thiosulfate ion were considered to be +6 as in sulfate and −2 as in sulfide for the central and terminal atoms, respectively. This view precluded the disproportionation reaction of thiosulfate into sulfate and sulfide as a redox mechanism for providing energy to bacteria under anaerobic conditions in sediments because there is no change in oxidation state for either S atom. However, XANES spectroscopy measurements have revealed that the charge densities of the sulfur atoms point towards +5 and −1 oxidation states for the central and terminal S atoms, respectively. This observation is consistent with the disproportionation of thiosulfate into sulfate and sulfide as a redox mechanism freeing up energy from microbial fermentation. Yet another interpretation suggests an oxidation state of +4 for the central S atom and 0 for the terminal atom and an unusually long 'full' S=S double bond between the two.
Formation
Thiosulfate ion is produced by the reaction of sulfite ion with elemental sulfur, and by incomplete oxidation of sulfides (e.g. pyrite oxidation). Sodium thiosulfate can be formed by disproportionation of sulfur dissolving in sodium hydroxide (similar to phosphorus).
Reactions
Thiosulfate ions reacts with acids to give sulfur dioxide and various sulfur rings:
This reaction may be used to generate sulfur colloids and demonstrate the Rayleigh scattering of light in physics. If white light is shone from below, blue light is seen from sideways and orange light from above, due to the same mechanisms that color the sky at midday and dusk.
Thiosulfate ions react with iodine to give tetrathionate ions:
This reaction is key for iodometry.
With bromine (X = Br) and chlorine (X = Cl), thiosulfate ions are oxidized to sulfate ions:
Reactions with metals and metal ions
Thiosulfate ion extensively forms diverse complexes with transition metals. This reactivity is related to its role in of silver-based photography.
Also reflecting its affinity for metals, thiosulfate ion rapidly corrodes metals in acidic conditions. Steel and stainless steel are particularly sensitive to pitting corrosion induced by thiosulfate ions. Molybdenum improves the resistance of stainless steel toward pitting (AISI 316L hMo). In alkaline aqueous conditions and medium temperature (60 °C), carbon steel and stainless steel (AISI 304L, 316L) are not attacked, even at high concentration of base (30%w KOH), thiosulfate ion (10%w) and in presence of fluoride ion (5%w KF).
In the era of silver-based photography, thiosulfate salts were consumed on a large scale as a "fixer" reagent. This application exploits thiosulfate ion's ability to dissolve silver halides. Sodium thiosulfate, commonly called hypo (from "hyposulfite"), was widely used in photography to fix black and white negatives and prints after the developing stage; modern "rapid" fixers use ammonium thiosulfate as a fixing salt because it acts three to four times faster.
Thiosulfate salts have been used to extract or leach gold and silver from their ores as a less toxic alternative to cyanide ion.
Biochemistry
The enzyme rhodanase (thiosulfate sulfurtransferase) catalyzes the detoxification of cyanide ion by thiosulfate ion by transforming them into thiocyanate ion and sulfite ion:
Sodium thiosulfate has been considered as an empirical treatment for cyanide poisoning, along with hydroxocobalamin. It is most effective in a pre-hospital setting, since immediate administration by emergency personnel is necessary to reverse rapid intracellular hypoxia caused by the inhibition of cellular respiration, at complex IV.
It activates thiosulfate sulfurtransferase (TST) in mitochondria. TST is associated with protection against obesity and type II (insulin resistant) diabetes.
Thiosulfate can also work as electron donor for growth of bacteria oxidizing sulfur, such as Chlorobium limicola forma thiosulfatophilum. These bacteria use electrons from thiosulfate (and other sources) and carbon from carbon dioxide to synthesize carbon compounds through reverse Krebs cycle.
Some bacteria can metabolise thiosulfates.
Minerals
Thiosulfate ion is a component of the very rare mineral sidpietersite . The presence of this anion in the mineral bazhenovite was disputed.
Nomenclature
Thiosulfate is an acceptable common name and used almost always.
The functional replacement IUPAC name is sulfurothioate; the systematic additive IUPAC name is trioxidosulfidosulfate(2−) or trioxido-1κ3O-disulfate(S—S)(2−).
Thiosulfate also refers to the esters of thiosulfuric acid, e.g. O,S-dimethyl thiosulfate . Such species are rare.
| Physical sciences | Sulfuric oxyanions | Chemistry |
8088857 | https://en.wikipedia.org/wiki/Cone%20snail | Cone snail | Cone snails, or cones, are highly venomous sea snails of the family Conidae.
Fossils of cone snails have been found from the Eocene to the Holocene epochs. Cone snail species have shells that are roughly conical in shape. Many species have colorful patterning on the shell surface. Cone snails are almost exclusively tropical in distribution.
All cone snails are venomous and capable of stinging. Cone snails use a modified radula tooth and a venom gland to attack and paralyze their prey before engulfing it. The tooth, which is likened to a dart or a harpoon, is barbed and can be extended some distance out from the head of the snail at the end of the proboscis.
Cone snail venoms are mainly peptide-based, and contain many different toxins that vary in their effects. The sting of several larger species of cone snails can be serious, and even fatal to humans. Cone snail venom also shows promise for medical use.
Distribution and habitat
There are over 900 different species of cone snails. Cone snails are typically found in warm tropical seas and oceans worldwide. Cone snails reach their greatest diversity in the Western Indo-Pacific region. While the majority of cone snails are found in warm tropical waters, some species have adapted to temperate/semi-tropical environments and are endemic to areas such as the Cape coast of South Africa, the Mediterranean, or the cool subtropical waters of southern California (Californiconus californicus).
Cone snails are found in all tropical and subtropical seas. They live on a variety of substrates, from the intertidal zone and deeper areas, to sand, rocks or coral reefs.
Shell
Cone snails have a large variety of shell colors and patterns, with local varieties and color forms of the same species often occurring. This variety in color and pattern has led to the creation of a large number of known synonyms and probable synonyms, making it difficult to give an exact taxonomic assignment for many snails in this genus. As of 2009, more than 3,200 different species names have been assigned, with an average of 16 new species names introduced each year.
The shells of cone snails vary in size and are conical in shape. The shell is whorled in the form of an inverted cone, with the anterior end being narrower. The protruding parts of the top of the whorls, that form the spire, are in the shape of another more flattened cone. The aperture is elongated and narrow with the sharp operculum being very small. The outer lip is simple, thin, and sharp, without a callus, and has a notched tip at the upper part. The columella is straight.
The larger species of cone snails can grow up to in length. The shells of cone snails are often brightly colored with a variety of patterns. Some species color patterns may be partially or completely hidden under an opaque layer of periostracum. In other species, the topmost shell layer is a thin periostracum, a transparent yellowish or brownish membrane.
Physiology and behavior
Cone snails are carnivorous. Their prey consists of marine worms, small fish, molluscs, and other cone snails. Cone snails are slow-moving, and use a venomous harpoon to disable faster-moving prey.
The osphradium in cone snails is more specialized than in other groups of gastropods. It is through this sensory modality that cone snails are able to sense their prey. The cone snails immobilize their prey using a modified, dartlike, barbed radular tooth, made of chitin, along with a venom gland containing neurotoxins.
Molecular phylogeny research has shown that preying on fish has evolved at least twice independently in cone snails. Some species appear to have also evolved prey mimicry, where they release chemicals that resemble the sex pheromones certain ragworms release during their short breeding season. The researchers hypothesize that these chemicals cause the prey to be more easily harpooned, but are still uncertain as to exactly how this occurs in the wild.
Harpoon
Cone snails use a harpoon-like structure called a radula tooth for predation. Radula teeth are modified teeth, primarily made of chitin and formed inside the mouth of the snail, in a structure known as the toxoglossan radula. Each specialized cone snail tooth is stored in the radula sac, except for the tooth that is in current use.
The radula tooth is hollow and barbed, and is attached to the tip of the radula in the radular sac, inside the snail's throat. When the snail detects a prey animal nearby, it extends a long flexible tube called a proboscis towards the prey. The radula tooth is loaded with venom from the venom bulb and, still attached to the radula, is fired from the proboscis into the prey by a powerful muscular contraction. The venom can paralyze smaller fish almost instantly. The snail then retracts the radula, drawing the subdued prey into the mouth. After the prey has been digested, the cone snail will regurgitate any indigestible material, such as spines and scales, along with the harpoon. There is always a radular tooth in the radular sac. A tooth may also be used in self-defense when the snail feels threatened.
The venom of cone snails contains hundreds of different compounds, and its exact composition varies widely from one species to another. The toxins in cone snail venom are referred to as conotoxins, and are composed of various peptides, each targeting a specific nerve channel or receptor. Some cone snail venoms also contain a pain-reducing toxin.
Relevance to humans
Dangers
Cone snails are prized for their brightly colored and patterned shells, which may tempt people to pick them up. This is risky, as the snail often fires its harpoon in self defense when disturbed. The harpoons of some of the larger species of cone snail can penetrate gloves or wetsuits.
The sting of many of the smallest cone species may be no worse than a bee or hornet sting, but the sting of a few of the larger tropical fish-eating species, such as Conus geographus, Conus tulipa and Conus striatus, can be fatal. Other dangerous species are Conus pennaceus, Conus textile, Conus aulicus, Conus magus and Conus marmoreus. According to Goldfrank's Toxicologic Emergencies, about 27 human deaths can be confidently attributed to cone snail envenomation, though the actual number is almost certainly much higher; some three dozen people are estimated to have died from geography cone envenomation alone.
Most of the cone snails that hunt worms are not a risk to humans, with the exception of larger species. One of the fish-eating species, the geography cone, Conus geographus, is also known colloquially as the "cigarette snail", a gallows humor exaggeration implying that, when stung by this creature, the victim will have only enough time to smoke a cigarette before dying.
Symptoms of a more serious cone snail sting include severe, localized pain, swelling, numbness and tingling, and vomiting. Symptoms can start immediately or can be delayed for days. Severe cases involve muscle paralysis, changes in vision and respiratory failure that can lead to death. If stung, one should seek medical attention as soon as possible.
Medical use
The appeal of conotoxins for creating pharmaceutical drugs is the precision and speed with which the chemicals act; many of the compounds target only a particular class of receptor. This means that they can reliably and quickly produce a particular effect on the body's systems without side effects; for example, almost instantly reducing heart rate or turning off the signaling of a single class of nerve, such as pain receptors.
Ziconotide, a pain reliever 1,000 times as powerful as morphine, was initially isolated from the venom of the magician cone snail, Conus magus. It was approved by the U.S. Food and Drug Administration in December 2004 under the name Prialt. Other drugs based on cone snail venom targeting Alzheimer's disease, Parkinson's disease, depression, and epilepsy are in clinical or preclinical trials.
Many peptides produced by the cone snails show prospects for being potent pharmaceuticals, such as AVC1, isolated from the Australian species, the Queen Victoria cone, Conus victoriae, and have been highly effective in treating postsurgical and neuropathic pain, even accelerating recovery from nerve injury.
Geography and tulip cone snails are known to secrete a type of insulin that paralyzes nearby fish by causing hypoglycaemic shock. They are the only two non-human animal species known to use insulin as a weapon. Cone snail insulin is capable of binding to human insulin receptors and researchers are studying its use as a potent fast-acting therapeutic insulin.
Shell collecting
The intricate color patterns of cone snails have made them one of the most popular species for shell collectors.
Conus gloriamaris, also known as "Glory of the Seas", one of the most famous and sought-after seashells in past centuries, with only a few specimens in private collections. The rarity of this species' shells led to high market prices for the objects, until the habitat of this cone snail was discovered, which decreased prices dramatically.
As jewelry
Naturally occurring, beach-worn cone shell tops can function as beads without any further modification. In Hawaii, these natural beads were traditionally collected from the beach drift to make puka shell jewelry. Since it is difficult to obtain enough naturally occurring cone snail tops, almost all modern puka shell jewelry uses cheaper imitations, cut from thin shells of other species of mollusk, or made of plastic.
Species
Until 2009 all species within the family Conidae were placed in one genus, Conus. Testing of the molecular phylogeny of the Conidae was first conducted by Christopher Meyer and Alan Kohn, and has continued, particularly with the advent of nuclear DNA testing.
In 2009, J.K. Tucker and M.J. Tenorio proposed a classification system consisting of three distinct families and 82 genera for living species of cone snails. This classification is based on shell morphology, radular differences, anatomy, physiology, and cladistics, with comparisons to molecular (DNA) studies. Published accounts of Conidae that use these new genera include J.K. Tucker & M.J. Tenorio (2009), and Bouchet et al. (2011). Tucker and Tenorio's proposed classification system for the cone shells and other clades of Conoidean gastropods is shown in Tucker & Tenorio cone snail taxonomy 2009.
Some experts, however, still prefer to use the traditional classification. For example, in the November 2011 version of the World Register of Marine Species, all species within the family Conidae were placed in the genus Conus. The binomial names of species in the 82 genera of living cone snails listed in Tucker & Tenorio 2009 were recognized by the World Register of Marine Species as "alternative representations". Debate within the scientific community regarding this issue has continued, and additional molecular phylogeny studies are being carried out in an attempt to clarify the issue.
In 2015, in the Journal of Molluscan Studies, Puillandre, Duda, Meyer, Olivera & Bouchet presented a new classification for the old genus Conus. Using 329 species, the authors carried out molecular phylogenetic analyses. The results suggested that the authors should place all cone snails in a single family, Conidae, containing four genera: Conus, Conasprella, Profundiconus and Californiconus. The authors group 85% of all known cone snail species under Conus. They recognize 57 subgenera within Conus, and 11 subgenera within the genus Conasprella.
Taxonomy
Afonsoconus Tucker & Tenorio, 2013
Africonus Petuch, 1975
Afroconus Petuch, 1975
Ammirales Schepman, 1913
Asperi Schepman, 1913
Asprella Schaufuss, 1869
Atlanticonus Petuch & Sargent, 2012
Attenuiconus Petuch, 2013
Austroconus Tucker & Tenorio, 2009
Bermudaconus Petuch, 2013
Brasiliconus Petuch, 2013
Calibanus da Motta, 1991
Cariboconus Petuch, 2003
Chelyconus Mörch, 1842
Cleobula 1930
Conasprelloides Tucker & Tenorio, 2009
Coronaxis Swainson, 1840
Cucullus Röding, 1798
Cylinder Montfort, 1810
Cylindrus Deshayes, 1824
Darioconus Iredale, 1930
Dauciconus Cotton, 1945
Dendroconus Swainson, 1840
Ductoconus da Motta, 1991
Embrikena Iredale, 1937
Endemoconus Iredale, 1931
Erythroconus da Motta, 1991
Eugeniconus da Motta, 1991
Floraconus Iredale, 1930
Gastridium Mödeer, 1793
Gladioconus Tucker & Tenorio, 2009
Gradiconus da Motta, 1991
Hermes Montfort, 1810
Heroconus da Motta, 1991
Isoconus Tucker & Tenorio, 2013
Kermasprella Powell, 1958
Ketyconus da Motta, 1991
Kioconus da Motta, 1991
Lautoconus Monterosato, 1923
Leporiconus Iredale, 1930
Leptoconus Swainson, 1840
Lilliconus Raybaudi Massilia, 1994
Lithoconus Mörch, 1852
Magelliconus da Motta, 1991
Mamiconus Cotton & Godfrey, 1932
Nitidoconus Tucker & Tenorio, 2013
Ongoconus da Motta, 1991
Phasmoconus Mörch, 1852
Pionoconus Mörch, 1852
Poremskiconus Petuch, 2013
Profundiconus Kuroda, 1956
Stephanoconus Mörch, 1852
Textilia Swainson, 1840
Tuliparia Swainson, 1840
Turriconus Shikama & Habe, 1968
Virgiconus Cotton, 1945
Virroconus Iredale, 1930
| Biology and health sciences | Gastropods | Animals |
246747 | https://en.wikipedia.org/wiki/Giant%20salamander | Giant salamander | The Cryptobranchidae (commonly known as giant salamanders) are a family of large salamanders that are fully aquatic. The family includes some of the largest living amphibians. They are native to China, Japan, and the eastern United States. Giant salamanders constitute one of two living families—the other being the Asiatic salamanders belonging to the family Hynobiidae—within the Cryptobranchoidea, one of two main divisions of living salamanders.
The largest species are in the genus Andrias, native to east Asia. The South China giant salamander (Andrias sligoi), can reach a length of . The Japanese giant salamander (Andrias japonicus) reaches up to in length, feeds at night on fish and crustaceans, and has been known to live for more than 50 years in captivity.
The hellbender (Cryptobranchus alleganiensis) inhabits the eastern United States and is the only member of the genus Cryptobranchus.
Taxonomy
The family name is from the Ancient Greek krypto ("hidden"), and branch ("gill"), which refer to how the members absorb oxygen through capillaries of their side-frills, which function as gills.
Clade Pancryptobrancha (Cryptobranchidae + Ukrainurus)
Genus †Ukrainurus Ukraine, Miocene
†Ukrainurus hypsognathus
Genus †Chunerpeton? China, Middle Jurassic
†Chunerpeton tianyiensis
Family Cryptobranchidae
Genus Cryptobranchus (hellbenders)
†Cryptobranchus saskatchewanensis? Ravenscrag Formation, Canada, Paleocene
Cryptobranchus alleganiensis (hellbender)
Genus Andrias (Asian giant salamanders; sometimes classified among the Cryptobranchus)
Andrias cheni (Qimen giant salamander)
Andrias davidianus (Chinese giant salamander) – (Simplified Chinese: ; pinyin: ) (may actually be a species complex of 5 different species)
Andrias sligoi (South China giant salamander)
Andrias japonicus (Japanese giant salamander) – ()
Andrias jiangxiensis (Jiangxi giant salamander)
†Andrias matthewi North America, Miocene
†Andrias scheuchzeri Europe, Oligocene-Pliocene
Genus †Eoscapherpeton Central Asia, Late Cretaceous (Cenomanian-Campanian)
Genus †Aviturus Mongolia, Paleocene
†Aviturus exsecratus
Genus †Ulanurus Mongolia, Paleocene
†Ulanurus fractus
Genus †Zaissanurus Kazakhstan, Oligocene
†Zaissanurus beliajevae
Phylogeny
The following phylogeny is based on Vasilyan et al. (2013):
The well-represented Cretaceous Eoscapherpeton was not phylogenetically placed. The enigmatic "Cryptobranchus" saskatchewanensis of Paleocene Canada may actually represent a stem-cryptobranchid.
Description
Cryptobranchids are large and predominantly nocturnal salamanders that can reach a length of , though most are considerably smaller today. Despite being aquatic, they are poor swimmers and mostly just walk on the bottom. Swimming by undulatory locomotion is generally used just for short distance-escapes to hiding places. The body is stout with large folds of skin along the flanks and a heavy, laterally compressed tail. These folds help increase the animals' surface area, allowing them to absorb more oxygen from the water as the adults lacks gills and have poorly developed lungs. Like in the majority of salamander species, there are four toes on the fore limbs and five on the hind limbs. They have paedomorphic traits, meaning their metamorphosis from the larval stage is incomplete, so they lack eyelids and the adults retain gill slits (open in the hellbender, closed in Andrias). Eyes are small and the eyesight bad.
Distribution and habitat
In Japan, their natural habitats are threatened by dam-building. Ramps and staircases have been added to some dams to allow them to move upstream to areas where they spawn.
Behavior
A Japanese giant salamander lived for 52 years in captivity.
Feeding
The Chinese giant salamander eats aquatic insects, fish, frogs, crabs, and shrimp. They hunt mainly at night. As they have poor eyesight, they use sensory nodes on their heads and bodies to detect minute changes in water pressure, enabling them to find their prey.
Reproduction
During mating season, the salamanders travel upstream, where the female lays two strings of over 200 eggs each. Lacking the stereotypical courtship behaviors found in other species, the male fertilizes the eggs externally by releasing his sperm onto them, and then guards them for at least three months, until they hatch. Tail fanning also occurs in order to increase the oxygen supply for the eggs. At this point, the larvae live off their noticeable stored fat until ready to hunt. Once ready, they hunt as a group rather than individually.
Scientists at Hiroshima City Asa Zoological Park in Japan have recently discovered the male salamander will spawn with more than one female in his den. Only large males can occupy and guard a den. They guard the den against other males and sexually inactive females. Those that are sexually active are welcomed. On occasion, the male "den master" will also allow a second male (smaller male salamanders, named "satellite males", who do not have their own den) into the den; the reason for this is unclear.
Fossil record
Extant species in the family Cryptobranchidae are the modern-day members of a lineage that extends back tens of millions of years. The oldest known fossils of cryptobranchoids are known from the Middle Jurassic of China. Chunerpeton from the Middle Jurassic of China has been suggested to represent the oldest known cryptobranchid. However, some studies have found it to be a more basal cryptobranchoid not more closely related to Cryptobranchidae than to Hynobiidae. The next oldest cryptobranchid is Eoscapherpeton, known from numerous Late Cretaceous deposits in Central Asia, which is suggested to represent a stem-group to modern cryptobranchids. Modern crown group representatives appear during the Paleocene.
As the fossil record for the Cryptobranchidae shows an Asian origin for the family, how these salamanders made it to the eastern US has been a point of scientific interest. Research has indicated a dispersal via land bridge, with waves of adaptive radiation seeming to have swept the Americas from north to south.
In 1726, the Swiss physician Johann Jakob Scheuchzer described a fossil as Homo diluvii testis (Latin: Evidence of a diluvian human), believing it to be the remains of a human being who drowned in the biblical flood. The Teylers Museum in Haarlem, Netherlands, bought the fossil in 1802, where it is still exhibited. In 1812, the fossil was examined by Georges Cuvier, who recognized that it was not human. After being identified as a salamander, it was renamed Salamandra scheuchzeri by Holl in 1831. The genus Andrias was coined six years later by Tschudi. In doing so, both the genus, Andrias (which means "image of man"), and the specific name, scheuchzeri, ended up honouring Scheuchzer and his beliefs. It and the extant A. davidianus cannot be mutually distinguished, and the latter, only described in 1871, is therefore sometimes considered a synonym of the former.
| Biology and health sciences | Salamanders and newts | Animals |
246806 | https://en.wikipedia.org/wiki/Ruminant | Ruminant | Ruminants are herbivorous grazing or browsing artiodactyls belonging to the suborder Ruminantia that are able to acquire nutrients from plant-based food by fermenting it in a specialized stomach prior to digestion, principally through microbial actions. The process, which takes place in the front part of the digestive system and therefore is called foregut fermentation, typically requires the fermented ingesta (known as cud) to be regurgitated and chewed again. The process of rechewing the cud to further break down plant matter and stimulate digestion is called rumination. The word "ruminant" comes from the Latin ruminare, which means "to chew over again".
The roughly 200 species of ruminants include both domestic and wild species. Ruminating mammals include cattle, all domesticated and wild bovines, goats, sheep, giraffes, deer, gazelles, and antelopes. It has also been suggested that notoungulates also relied on rumination, as opposed to other atlantogenatans that rely on the more typical hindgut fermentation, though this is not entirely certain.
Ruminants represent the most diverse group of living ungulates. The suborder Ruminantia includes six different families: Tragulidae, Giraffidae, Antilocapridae, Cervidae, Moschidae, and Bovidae.
Taxonomy and evolution
The first fossil ruminants appeared in the Early Eocene and were small, likely omnivorous, forest-dwellers. Artiodactyls with cranial appendages first occur in the early Miocene.
Phylogeny
Ruminantia is a crown group of ruminants within the order Artiodactyla, cladistically defined by Spaulding et al. as "the least inclusive clade that includes Bos taurus (cow) and Tragulus napu (mouse deer)". Ruminantiamorpha is a higher-level clade of artiodactyls, cladistically defined by Spaulding et al. as "Ruminantia plus all extinct taxa more closely related to extant members of Ruminantia than to any other living species." This is a stem-based definition for Ruminantiamorpha, and is more inclusive than the crown group Ruminantia. As a crown group, Ruminantia only includes the last common ancestor of all extant (living) ruminants and their descendants (living or extinct), whereas Ruminantiamorpha, as a stem group, also includes more basal extinct ruminant ancestors that are more closely related to living ruminants than to other members of Artiodactyla. When considering only living taxa (neontology), this makes Ruminantiamorpha and Ruminantia synonymous, and only Ruminantia is used. Thus, Ruminantiamorpha is only used in the context of paleontology. Accordingly, Spaulding grouped some genera of the extinct family Anthracotheriidae within Ruminantiamorpha (but not in Ruminantia), but placed others within Ruminantiamorpha's sister clade, Cetancodontamorpha.
Ruminantia's placement within Artiodactyla can be represented in the following cladogram:
Within Ruminantia, the Tragulidae (mouse deer) are considered the most basal family, with the remaining ruminants classified as belonging to the infraorder Pecora. Until the beginning of the 21st century it was understood that the family Moschidae (musk deer) was sister to Cervidae. However, a 2003 phylogenetic study by Alexandre Hassanin (of National Museum of Natural History, France) and colleagues, based on mitochondrial and nuclear analyses, revealed that Moschidae and Bovidae form a clade sister to Cervidae. According to the study, Cervidae diverged from the Bovidae-Moschidae clade 27 to 28 million years ago. The following cladogram is based on a large-scale genome ruminant genome sequence study from 2019:
Classification
Order Artiodactyla
Suborder Tylopoda: camels and llamas, 7 living species in 3 genera
Suborder Suina: pigs and peccaries
Suborder Cetruminantia: ruminants, whales and hippos
unranked Ruminantia
Infraorder Tragulina (paraphyletic)
Family †Leptomerycidae
Family †Hypertragulidae
Family †Praetragulidae
Family †Gelocidae
Family †Bachitheriidae
Family Tragulidae: chevrotains, 6 living species in 4 genera
Family †Archaeomerycidae
Family †Lophiomerycidae
Infraorder Pecora
Family Cervidae: deer and moose, 49 living species in 16 genera
Family †Palaeomerycidae
Family †Dromomerycidae
Family †Hoplitomerycidae
Family †Climacoceratidae
Family Giraffidae: giraffe and okapi, 2 living species in 2 genera
Family Antilocapridae: pronghorn, one living species in one genus
Family Moschidae: musk deer, 4 living species in one genus
Family Bovidae: cattle, goats, sheep, and antelope, 143 living species in 53 genera
Digestive system of ruminants
Hofmann and Stewart divided ruminants into three major categories based on their feed type and feeding habits: concentrate selectors, intermediate types, and grass/roughage eaters, with the assumption that feeding habits in ruminants cause morphological differences in their digestive systems, including salivary glands, rumen size, and rumen papillae. However, Woodall found that there is little correlation between the fiber content of a ruminant's diet and morphological characteristics, meaning that the categorical divisions of ruminants by Hofmann and Stewart warrant further research.
Also, some mammals are pseudoruminants, which have a three-compartment stomach instead of four like ruminants. The Hippopotamidae (comprising hippopotamuses) are well-known examples. Pseudoruminants, like traditional ruminants, are foregut fermentors and most ruminate or chew cud. However, their anatomy and method of digestion differs significantly from that of a four-chambered ruminant.
Monogastric herbivores, such as rhinoceroses, horses, guinea pigs, and rabbits, are not ruminants, as they have a simple single-chambered stomach. Being hindgut fermenters, these animals ferment cellulose in an enlarged cecum. In smaller hindgut fermenters of the order Lagomorpha (rabbits, hares, and pikas), and Caviomorph rodents (Guinea pigs, capybaras, etc.), material from the cecum is formed into cecotropes, passed through the large intestine, expelled and subsequently reingested to absorb nutrients in the cecotropes.
The primary difference between ruminants and nonruminants is that ruminants' stomachs have four compartments:
rumen—primary site of microbial fermentation
reticulum
omasum—receives chewed cud, and absorbs volatile fatty acids
abomasum—true stomach
The first two chambers are the rumen and the reticulum. These two compartments make up the fermentation vat and are the major site of microbial activity. Fermentation is crucial to digestion because it breaks down complex carbohydrates, such as cellulose, and enables the animal to use them. Microbes function best in a warm, moist, anaerobic environment with a temperature range of and a pH between 6.0 and 6.4. Without the help of microbes, ruminants would not be able to use nutrients from forages. The food is mixed with saliva and separates into layers of solid and liquid material. Solids clump together to form the cud or bolus.
The cud is then regurgitated and chewed to completely mix it with saliva and to break down the particle size. Smaller particle size allows for increased nutrient absorption. Fiber, especially cellulose and hemicellulose, is primarily broken down in these chambers by microbes (mostly bacteria, as well as some protozoa, fungi, and yeast) into the three volatile fatty acids (VFAs): acetic acid, propionic acid, and butyric acid. Protein and nonstructural carbohydrate (pectin, sugars, and starches) are also fermented. Saliva is very important because it provides liquid for the microbial population, recirculates nitrogen and minerals, and acts as a buffer for the rumen pH. The type of feed the animal consumes affects the amount of saliva that is produced.
Though the rumen and reticulum have different names, they have very similar tissue layers and textures, making it difficult to visually separate them. They also perform similar tasks. Together, these chambers are called the reticulorumen. The degraded digesta, which is now in the lower liquid part of the reticulorumen, then passes into the next chamber, the omasum. This chamber controls what is able to pass into the abomasum. It keeps the particle size as small as possible in order to pass into the abomasum. The omasum also absorbs volatile fatty acids and ammonia.
After this, the digesta is moved to the true stomach, the abomasum. This is the gastric compartment of the ruminant stomach. The abomasum is the direct equivalent of the monogastric stomach, and digesta is digested here in much the same way. This compartment releases acids and enzymes that further digest the material passing through. This is also where the ruminant digests the microbes produced in the rumen. Digesta is finally moved into the small intestine, where the digestion and absorption of nutrients occurs. The small intestine is the main site of nutrient absorption. The surface area of the digesta is greatly increased here because of the villi that are in the small intestine. This increased surface area allows for greater nutrient absorption. Microbes produced in the reticulorumen are also digested in the small intestine. After the small intestine is the large intestine. The major roles here are breaking down mainly fiber by fermentation with microbes, absorption of water (ions and minerals) and other fermented products, and also expelling waste. Fermentation continues in the large intestine in the same way as in the reticulorumen.
Only small amounts of glucose are absorbed from dietary carbohydrates. Most dietary carbohydrates are fermented into VFAs in the rumen. The glucose needed as energy for the brain and for lactose and milk fat in milk production, as well as other uses, comes from nonsugar sources, such as the VFA propionate, glycerol, lactate, and protein. The VFA propionate is used for around 70% of the glucose and glycogen produced and protein for another 20% (50% under starvation conditions).
Abundance, distribution, and domestication
Wild ruminants number at least 75 million and are native to all continents except Antarctica and Australia. Nearly 90% of all species are found in Eurasia and Africa. Species inhabit a wide range of climates (from tropic to arctic) and habitats (from open plains to forests).
The population of domestic ruminants is greater than 3.5 billion, with cattle, sheep, and goats accounting for about 95% of the total population. Goats were domesticated in the Near East circa 8000 BC. Most other species were domesticated by 2500 BC., either in the Near East or southern Asia.
Ruminant physiology
Ruminating animals have various physiological features that enable them to survive in nature. One feature of ruminants is their continuously growing teeth. During grazing, the silica content in forage causes abrasion of the teeth. This is compensated for by continuous tooth growth throughout the ruminant's life, as opposed to humans or other nonruminants, whose teeth stop growing after a particular age. Most ruminants do not have upper incisors; instead, they have a thick dental pad to thoroughly chew plant-based food. Another feature of ruminants is the large ruminal storage capacity that gives them the ability to consume feed rapidly and complete the chewing process later. This is known as rumination, which consists of the regurgitation of feed, rechewing, resalivation, and reswallowing. Rumination reduces particle size, which enhances microbial function and allows the digesta to pass more easily through the digestive tract.
Unlike camelids, ruminants copulate in a standing position and are not Induced ovulators.
Rumen microbiology
Vertebrates lack the ability to hydrolyse the beta [1–4] glycosidic bond of plant cellulose due to the lack of the enzyme cellulase. Thus, ruminants completely depend on the microbial flora, present in the rumen or hindgut, to digest cellulose. Digestion of food in the rumen is primarily carried out by the rumen microflora, which contains dense populations of several species of bacteria, protozoa, sometimes yeasts and other fungi – 1 ml of rumen is estimated to contain 10–50 billion bacteria and 1 million protozoa, as well as several yeasts and fungi.
Since the environment inside a rumen is anaerobic, most of these microbial species are obligate or facultative anaerobes that can decompose complex plant material, such as cellulose, hemicellulose, starch, and proteins. The hydrolysis of cellulose results in sugars, which are further fermented to acetate, lactate, propionate, butyrate, carbon dioxide, and methane.
As bacteria conduct fermentation in the rumen, they consume about 10% of the carbon, 60% of the phosphorus, and 80% of the nitrogen that the ruminant ingests. To reclaim these nutrients, the ruminant then digests the bacteria in the abomasum. The enzyme lysozyme has adapted to facilitate digestion of bacteria in the ruminant abomasum. Pancreatic ribonuclease also degrades bacterial RNA in the ruminant small intestine as a source of nitrogen.
During grazing, ruminants produce large amounts of saliva – estimates range from 100 to 150 litres of saliva per day for a cow. The role of saliva is to provide ample fluid for rumen fermentation and to act as a buffering agent. Rumen fermentation produces large amounts of organic acids, thus maintaining the appropriate pH of rumen fluids is a critical factor in rumen fermentation. After digesta passes through the rumen, the omasum absorbs excess fluid so that digestive enzymes and acid in the abomasum are not diluted.
Tannin toxicity in ruminant animals
Tannins are phenolic compounds that are commonly found in plants. Found in the leaf, bud, seed, root, and stem tissues, tannins are widely distributed in many different species of plants. Tannins are separated into two classes: hydrolysable tannins and condensed tannins. Depending on their concentration and nature, either class can have adverse or beneficial effects. Tannins can be beneficial, having been shown to increase milk production, wool growth, ovulation rate, and lambing percentage, as well as reducing bloat risk and reducing internal parasite burdens.
Tannins can be toxic to ruminants, in that they precipitate proteins, making them unavailable for digestion, and they inhibit the absorption of nutrients by reducing the populations of proteolytic rumen bacteria. Very high levels of tannin intake can produce toxicity that can even cause death. Animals that normally consume tannin-rich plants can develop defensive mechanisms against tannins, such as the strategic deployment of lipids and extracellular polysaccharides that have a high affinity to binding to tannins. Some ruminants (goats, deer, elk, moose) are able to consume food high in tannins (leaves, twigs, bark) due to the presence in their saliva of tannin-binding proteins.
Religious importance
The Law of Moses in the Bible allowed the eating of some mammals that had cloven hooves (i.e. members of the order Artiodactyla) and "that chew the cud", a stipulation preserved to this day in Jewish dietary laws.
Other uses
The verb 'to ruminate' has been extended metaphorically to mean to ponder thoughtfully or to meditate on some topic. Similarly, ideas may be 'chewed on' or 'digested'. 'Chew the (one's) cud' is to reflect or meditate. In psychology, "rumination" refers to a pattern of thinking, and is unrelated to digestive physiology.
Ruminants and climate change
Methane is produced by a type of archaea, called methanogens, as described above within the rumen, and this methane is released to the atmosphere. The rumen is the major site of methane production in ruminants. Methane is a strong greenhouse gas with a global warming potential of 86 compared to CO2 over a 20-year period.
As a by-product of consuming cellulose, cattle belch out methane, there-by returning that carbon sequestered by plants back into the atmosphere. After about 10 to 12 years, that methane is broken down and converted back to . Once converted to , plants can again perform photosynthesis and fix that carbon back into cellulose. From here, cattle can eat the plants and the cycle begins once again. In essence, the methane belched from cattle is not adding new carbon to the atmosphere. Rather it is part of the natural cycling of carbon through the biogenic carbon cycle.
In 2010, enteric fermentation accounted for 43% of the total greenhouse gas emissions from all agricultural activity in the world, 26% of the total greenhouse gas emissions from agricultural activity in the U.S., and 22% of the total U.S. methane emissions. The meat from domestically raised ruminants has a higher carbon equivalent footprint than other meats or vegetarian sources of protein based on a global meta-analysis of lifecycle assessment studies. Methane production by meat animals, principally ruminants, is estimated 15–20% global production of methane, unless the animals were hunted in the wild. The current U.S. domestic beef and dairy cattle population is around 90 million head, approximately 50% higher than the peak wild population of American bison of 60 million head in the 1700s, which primarily roamed the part of North America that now makes up the United States.
| Biology and health sciences | Artiodactyla | null |
246828 | https://en.wikipedia.org/wiki/Forage | Forage | Forage is a plant material (mainly plant leaves and stems) eaten by grazing livestock. Historically, the term forage has meant only plants eaten by the animals directly as pasture, crop residue, or immature cereal crops, but it is also used more loosely to include similar plants cut for fodder and carried to the animals, especially as hay or silage.
While the term forage has a broad definition, the term forage crop is used to define crops, annual or biennial, which are grown to be utilized by grazing or harvesting as a whole crop.
Common forages
Grasses
Grass forages include:
Agrostis spp. – bentgrasses
Agrostis capillaris – common bentgrass
Agrostis stolonifera – creeping bentgrass
Andropogon hallii – sand bluestem
Arrhenatherum elatius – false oat-grass
Bothriochloa bladhii – Australian bluestem
Bothriochloa pertusa – hurricane grass
Brachiaria decumbens – Surinam grass
Brachiaria humidicola – koronivia grass
Bromus spp. – bromegrasses
Cenchrus ciliaris – buffelgrass
Chloris gayana – Rhodes grass
Cynodon dactylon – bermudagrass
Dactylis glomerata – orchard grass
Echinochloa pyramidalis – antelope grass
Entolasia imbricata – bungoma grass
Festuca spp. – fescues
Festuca arundinacea – tall fescue
Festuca pratensis – meadow fescue
Festuca rubra – red fescue
Heteropogon contortus – black spear grass
Hymenachne amplexicaulis – West Indian marsh grass
Hyparrhenia rufa – jaragua
Leersia hexandra – southern cutgrass
Lolium spp. – ryegrasses
Lolium multiflorum – Italian ryegrass
Lolium perenne – perennial ryegrass
Megathyrsus maximus – Guinea grass
Melinis minutiflora – molasses grass
Paspalum conjugatum – carabao grass
Paspalum dilatatum – dallisgrass
Phalaris arundinacea – reed canarygrass
Phleum pratense – timothy
Poa spp. – bluegrasses, meadow-grasses
Poa arachnifera – Texas bluegrass
Poa pratensis – Kentucky bluegrass
Poa trivialis – rough bluegrass
Setaria sphacelata – African bristlegrass
Themeda triandra – kangaroo grass
Thinopyrum intermedium – intermediate wheatgrass
Herbaceous legumes
Herbaceous legume forages include:
Arachis pintoi – pinto peanut
Astragalus cicer – cicer milkvetch
Chamaecrista rotundifolia – roundleaf sensitive pea
Clitoria ternatea – butterfly-pea
Kummerowia – annual lespedezas
Kummerowia stipulacea – Korean clover, Korean lespedeza
Kummerowia striata – Japanese clover, common lespedeza
Lotus corniculatus – bird's-foot trefoil
Macroptilium atropurpureum – purple bush-bean
Macroptilium bracteatum – burgundy bean
Medicago spp. – medics
Medicago sativa – alfalfa, lucerne
Medicago truncatula – barrel medic
Melilotus spp. – sweetclovers
Neonotonia wightii – perennial soybean
Onobrychis viciifolia – common sainfoin
Stylosanthes spp. – stylo
Stylosanthes humilis – Townsville stylo
Stylosanthes scabra – shrubby stylo
Trifolium spp. – clovers
Trifolium hybridum – alsike clover
Trifolium incarnatum – crimson clover
Trifolium pratense – red clover
Trifolium repens – white clover
Vicia spp. – vetches
Vicia articulata – oneflower vetch
Vicia ervilia – bitter vetch
Vicia narbonensis – narbon vetch
Vicia sativa – common vetch, tare
Vicia villosa – hairy vetch
Vigna parkeri – creeping vigna
Tree legumes
Tree legume forages include:
Acacia aneura – mulga
Albizia spp. – silk trees
Albizia canescens – Belmont siris
Albizia lebbeck – lebbeck
Enterolobium cyclocarpum – earpodtree
Leucaena leucocephala – leadtree
Silage
Silage may be composed by the following:
Alfalfa
Maize (corn)
Grass-legume mix
Sorghums
Oats
Aquatic feeds
Lemna minor – Duckweed
Pistia stratiotes – Water lecttuce
Eichhornia crassipes – Water hyacinth
Salvinia molesta – fern
Ipomoea aquatica – Water spinach
Crop residue
Crop residues used as forage include:
Sorghum
Sweet potato vines
Corn or soybean
Fruit tree by-products stover
Other
Raphanus sativus var. longipinnatus – Daikon radish/"forage radish"
| Technology | Animal husbandry | null |
246940 | https://en.wikipedia.org/wiki/Malvaceae | Malvaceae | Malvaceae (), or the mallows, is a family of flowering plants estimated to contain 244 genera with 4225 known species. Well-known members of economic importance include okra, cotton, cacao, roselle and durian. There are also some genera containing familiar ornamentals, such as Alcea (hollyhock), Malva (mallow), and Tilia (lime or linden tree). The genera with the largest numbers of species include Hibiscus (434 species), Pavonia (291 species), Sida (275 species), Ayenia (216 species), Dombeya (197 species), and Sterculia (181 species).
Taxonomy and nomenclature
The circumscription of the Malvaceae is controversial. The traditional Malvaceae sensu stricto comprise a very homogeneous and cladistically monophyletic group. Another major circumscription, Malvaceae sensu lato, has been more recently defined on the basis that genetics studies have shown the commonly recognised families Bombacaceae, Tiliaceae, and Sterculiaceae, which have always been considered closely allied to Malvaceae s.s., are not monophyletic groups. The Malvaceae can be expanded to include all of these families so as to compose a monophyletic group. Adopting this circumscription, the Malvaceae incorporate a much larger number of genera.
Subfamilies
This article is based on the second circumscription, as presented by the Angiosperm Phylogeny Website. The Malvaceae s.l. (hereafter simply "Malvaceae") comprise nine subfamilies. A tentative cladogram of the family is shown below. The diamond denotes a poorly supported branching (<80%).
Until recently, relationships between these subfamilies were either poorly supported or almost completely obscure. Continuing disagreements focused primarily on the correct circumscription of these subfamilies, including the preservation of the family Bombacaceae. A study published in 2021 presented a fully resolved phylogenetic framework for Malvaceae s.l. using genomic data for all nine subfamilies.
Regarding the traditional Malvaceae s.s., the subfamily Malvoideae approximately corresponds to that group.
Genera
245 genera are currently accepted.
Synapomorphies
The relationships between the "core Malvales" families used to be defined on the basis of shared "malvean affinities". These included the presence of malvoid teeth, stems with mucilage canals, and stratified wedge-shaped phloem. These affinities were problematic because they were not always shared within the core families. Later studies revealed more unambiguous synapomorphies within Malvaceae s.l.. Synapomorphies identified within Malvaceae s.l. include the presence of tile cells, trichomatous nectaries, and an inflorescence structure called a bicolor unit. Tile cells consist of vertically positioned cells interspersed between and dimensionally similar to procumbent ray cells. Evidence of Malvean wood fossils has confirmed their evolutionary link in Malvaceae s.l., as well as explained their diverse structures. Flowers of Malvaceae s.l. exhibit nectaries consisting of densely arranged multicellular hairs resembling trichomes. In most of Malvaceae s.l., these trichomatous nectaries are located on the inner surface of the sepals, but flowers of the subfamily Tiliodeae also have present nectaries on the petals. Malvean flowers also share a unifying structure known as a bicolor unit, named for its initial discovery in the flowers of Theobroma bicolor. The bicolor unit consists of an ordered inflorescence with determinate cymose structures. The inflorescence can branch off the main axis, creating separate orders of the flowers, with the main axis developing first. Bracts on the peduncle subtend axillary buds that become these lateral stalks. One bract within this whorl is a sterile bract. The bicolor unit is a variable structure in complexity, but the presence of fertile and sterile bracts is a salient characteristic.
Names
The English common name 'mallow' (also applied to other members of Malvaceae) comes from Latin malva (also the source for the English word "mauve"). Malva itself was ultimately derived from the word for the plant in ancient Mediterranean languages. Cognates of the word include Ancient Greek () or (), Modern Greek (), modern () and modern ().
Description
Most species are herbaceous plants or shrubs, but some are trees or lianas.
Leaves and stems
Leaves are generally alternate, often palmately lobed or compound and palmately veined. The margin may be entire, but when dentate, a vein ends at the tip of each tooth (malvoid teeth). Stipules are present. The stems contain mucous canals and often also mucous cavities. Hairs are common, and are most typically stellate. Stems of Bombacoideae are often covered in thick prickles.
Flowers
The flowers are commonly borne in definite or indefinite axillary inflorescences, which are often reduced to a single flower, but may also be cauliflorous, oppositifolious, or terminal. They often bear supernumerary bracts in the structure of a bicolor unit. They can be unisexual or bisexual, and are generally actinomorphic, often associated with conspicuous bracts, forming an epicalyx. They generally have five valvate sepals, most frequently basally connate, with five imbricate petals. The stamens are five to numerous, and connate at least at their bases, but often forming a tube around the pistils. The pistils are composed of two to many connate carpels. The ovary is superior, with axial placentation, with capitate or lobed stigma.
The flowers have nectaries made of many tightly packed glandular hairs, usually positioned on the sepals.
Fruits
The fruits are most often loculicidal capsules, schizocarps or nuts.
Pollination
Self-pollination is often avoided by means of protandry. Most species are entomophilous (pollinated by insects). Bees from the tribe Emphorini of the Apidae (including Ptilothrix, Diadasia, and Melitoma) are known to specialize on the plants.
Importance
A number of species are pests in agriculture, including Abutilon theophrasti and Modiola caroliniana, and others that are garden escapees. Cotton (four species of Gossypium), kenaf (Hibiscus cannabinus), cacao (Theobroma cacao), kola nut (Cola spp.), and okra (Abelmoschus esculentus) are important agricultural crops. The fruit and leaves of baobabs are edible, as is the fruit of the durian. A number of species, including Hibiscus syriacus, Hibiscus rosa-sinensis and Alcea rosea are garden plants.
| Biology and health sciences | Others | null |
247151 | https://en.wikipedia.org/wiki/Mimicry | Mimicry | In evolutionary biology, mimicry is an evolved resemblance between an organism and another object, often an organism of another species. Mimicry may evolve between different species, or between individuals of the same species. In the simplest case, as in Batesian mimicry, a mimic resembles a model, so as to deceive a dupe, all three being of different species. A Batesian mimic, such as a hoverfly, is harmless, while its model, such as a wasp, is harmful, and is avoided by the dupe, such as an insect-eating bird. Birds hunt by sight, so the mimicry in that case is visual, but in other cases mimicry may make use of any of the senses. Most types of mimicry, including Batesian, are deceptive, as the mimics are not harmful, but Müllerian mimicry, where different harmful species resemble each other, is honest, as when species of wasps and of bees all have genuinely aposematic warning coloration. More complex types may be bipolar, involving only two species, such as when the model and the dupe are the same; this occurs for example in aggressive mimicry, where a predator in wolf-in-sheep's-clothing style resembles its prey, allowing it to hunt undetected. Mimicry is not limited to animals; in Pouyannian mimicry, an orchid flower is the mimic, resembling a female bee, its model; the dupe is the male bee of the same species, which tries to copulate with the flower, enabling it to transfer pollen, so the mimicry is again bipolar. In automimicry, another bipolar system, model and mimic are the same, as when blue lycaenid butterflies have 'tails' or eyespots on their wings that mimic their own heads, misdirecting predator dupes to strike harmlessly. Many other types of mimicry exist.
Etymology
Use of the word mimicry dates to 1637. It derives from the Greek term mimetikos, "imitative", in turn from mimetos, the verbal adjective of mimeisthai, "to imitate". "Mimicry" was first used in zoology by the English entomologists William Kirby and William Spence in 1823. Originally used to describe people, "mimetic" was used in zoology from 1851.
History
Ancient
Aristotle wrote in his History of Animals that partridges use a deceptive distraction display to lure predators away from their flightless young:
The behaviour is recognised as a form of mimicry by biologists.
19th century
In 1823, Kirby and Spence, in their book An Introduction to Entomology, used the term "mimicry" informally to depict the way that the structure and coloration of some insects resembled objects in their environments:
The English naturalist Henry Walter Bates worked for several years on butterflies in the Amazon rainforest. Returning home, he described multiple forms of mimicry in an 1862 paper at the Linnean Society in London, and then in his 1863 book The Naturalist on the River Amazons. The term "Batesian mimicry" has since been used in his honour, its usage becoming restricted to the situation in which a harmless mimic gains protection from its predators by resembling a distasteful model. Among the observations in Bates's 1862 paper is the statement:
The German naturalist Fritz Müller also spent many years studying butterflies in the Amazon rainforest. He first published a journal article on mimicry in German in 1878, followed in 1879 by a paper to the Entomological Society of London (translated and presented by Ralph Meldola). He described a situation where different species were each unpalatable to predators, and shared similar, genuine, warning signals. Bates found it hard to explain why this should be so, asking why they should need to mimic each other if both were harmful and could warn off predators on their own. Müller put forward the first mathematical model of mimicry for this phenomenon: if a common predator confuses the two species, individuals in both those species are more likely to survive, as fewer individuals of either species are killed by the predator. The term Müllerian mimicry, named in his honour, has since been used for this mutualistic form of mimicry.
Müller wrote that
Overview
Evolved resemblance
Mimicry is an evolved resemblance between an organism and another object, often an organism of another species. Mimicry may evolve between different species, or between individuals of the same species. Often, mimicry functions to protect from predators. Mimicry systems have three basic roles: a mimic, a model, and a dupe. When these correspond to three separate species, the system is called disjunct; when the roles are taken by just two species, the system is called bipolar. Mimicry evolves if a dupe (such as a predator) perceives a mimic (such as a palatable prey) as a model (the organism it resembles), and is deceived to change its behaviour to the mimic's selective advantage. The resemblances can be via any sensory modality, including any combination of visual, acoustic, chemical, tactile, or electric. Mimicry may be to the advantage of both organisms that share a resemblance, in which case it is mutualistic; or it can be to the detriment of one, making it parasitic or competitive. The evolutionary convergence between groups is driven by the selective action of a dupe. Birds, for example, use sight to identify palatable insects, whilst avoiding noxious ones. Over time, palatable insects may evolve to resemble noxious ones, making them mimics and the noxious ones models. Models do not have to be more abundant than mimics. In the case of mutualism, each model is also a mimic; all such species can be called "co-mimics". Many harmless species such as hoverflies are Batesian mimics of strongly defended species such as wasps, while many such well-defended species form Müllerian mimicry rings of co-mimics. In the evolution of wasp-like appearance, it has been argued that insects evolve to masquerade wasps since predatory wasps do not attack each other, and that this mimetic resemblance has had the useful side-effect of deterring vertebrate predators.
Mimicry can result in an evolutionary arms race if mimicry negatively affects the model, in which case the model can evolve a different appearance from the mimic.p161 Mimics may have different models for different life cycle stages, or they may be polymorphic, with different individuals imitating different models, as occurs in Heliconius butterflies. Models tend to be relatively closely related to their mimics, but mimicry can be of vastly different species, for example when spiders mimic ants. Most known mimics are insects, though many other examples including vertebrates, plants, and fungi exist.
Evolutionary explanations
It is widely accepted that mimicry evolves as a positive adaptation. The lepidopterist and novelist Vladimir Nabokov however argued that although natural selection might stabilize a "mimic" form, it would not be necessary to create it. The most widely accepted model used to explain the evolution of mimicry in butterflies is the two-step hypothesis. The first step involves mutation in modifier genes that regulate a complex cluster of linked genes that cause large changes in morphology. The second step consists of selections on genes with smaller phenotypic effects, creating an increasingly close resemblance. This model is supported by empirical evidence that suggests that a few single point mutations cause large phenotypic effects, while numerous others produce smaller effects. Some regulatory elements collaborate to form a supergene for the development of butterfly color patterns. The model is supported by computational simulations of population genetics. The Batesian mimicry in Papilio polytes is controlled by the doublesex gene.
Some mimicry is imperfect. Natural selection drives mimicry only far enough to deceive predators. For example, when predators avoid a mimic that imperfectly resembles a coral snake, the mimic is sufficiently protected.
Convergent evolution is an alternative explanation for why coral reef fish have come to resemble each other; the same applies to benthic marine invertebrates such as sponges and nudibranchs.
Living and non-living models
In its broadest definition, mimicry can include non-living models. The specific terms masquerade and mimesis are sometimes used when the models are inanimate, and the mimicry's purpose is crypsis. For example, animals such as flower mantises, planthoppers, comma and geometer moth caterpillars resemble twigs, bark, leaves, bird droppings or flowers. In addition, predators may make use of resemblance to harmless objects in aggressive masquerade, to enable them to approach prey. This wolf in sheep's clothing strategy differs from the more specific resemblance to the prey in aggressive mimicry, where the prey is both model and dupe.
Many animals bear eyespots, which are hypothesized to resemble the eyes of larger animals. They may not resemble any specific organism's eyes, and whether or not animals respond to them as eyes is also unclear. The model is usually another species, except in automimicry, where members of the species mimic other members, or other parts of their own bodies, and in inter-sexual mimicry, where members of one sex mimic members of the other.
Types
Many types of mimicry have been described. An overview of each follows, highlighting the similarities and differences between the various forms. Classification is often based on function with respect to the mimic (e.g., avoiding harm). Some cases may belong to more than one class, e.g., automimicry and aggressive mimicry are not mutually exclusive, as one describes the species relationship between model and mimic, while the other describes the function for the mimic (obtaining food). The terminology used has been debated, as classifications have differed or overlapped; attempts to clarify definitions have led to the partial replacement of old terms with new ones.
Defensive
Mimicry is defensive or protective when organisms are able to avoid harmful encounters by deceiving enemies into treating them as something else.
Batesian
In Batesian mimicry, the mimic resembles the model, but does not have the attribute that makes it unprofitable to predators (e.g., unpalatability, or the ability to sting). In other words, a Batesian mimic is a sheep in wolf's clothing. Mimics are less likely to be found out (for example by predators) when in low proportion to their model. Such negative frequency-dependent selection applies in most forms of mimicry. Specifically, Batesian mimicry can only be maintained if the harm caused to the predator by eating a model outweighs the benefit of eating a mimic. The nature of learning is weighted in favor of the mimics, for a predator that has a bad first experience with a model tends to avoid anything that looks like it for a long time, and does not re-sample soon to see whether the initial experience was a false negative. However, if mimics become more abundant than models, then the probability of a young predator having a first experience with a mimic increases. Batesian systems are therefore most likely to be stable where the model is more abundant than the mimic.
There are many Batesian mimics among butterflies and moths. Consul fabius and Eresia eunice imitate unpalatable Heliconius butterflies such as H. ismenius. Limenitis arthemis imitate the poisonous pipevine swallowtail (Battus philenor). Several palatable moths produce ultrasonic click calls to mimic unpalatable tiger moths. Octopuses of the genus Thaumoctopus (the mimic octopus) are able to intentionally alter their body shape and coloration to resemble dangerous sea snakes or lionfish. In the Amazon, the helmeted woodpecker (Dryocopus galeatus), a rare species which lives in the Atlantic Forest of Brazil, Paraguay, and Argentina, has a similar red crest, black back, and barred underside to two larger woodpeckers: Dryocopus lineatus and Campephilus robustus. This mimicry reduces attacks on D. galeatus. Batesian mimicry occurs in the plant kingdom, where the chameleon vine adapts its leaf shape and colour to match that of the plant it is climbing.
Müllerian
In Müllerian mimicry, two or more species have similar warning or aposematic signals and both share genuine anti-predation attributes (e.g. being unpalatable), as first described in Heliconius butterflies. This type of mimicry is unique in several respects. Firstly, both the mimic and the model benefit from the interaction, which could thus be classified as mutualism. The signal receiver also benefits by this system, despite being deceived about species identity, as it is able to generalize the pattern to potentially harmful encounters. The distinction between mimic and model that is clear in Batesian mimicry is also blurred. Where one species is scarce and another abundant, the rare species can be said to be the mimic. When both are present in similar numbers, however, it makes more sense to speak of each as a co-mimic than of distinct 'mimic' and 'model' species, as their warning signals tend to converge. Also, the mimetic species may exist on a continuum from harmless to highly noxious, so Batesian mimicry grades smoothly into Müllerian convergence.
Emsleyan/Mertensian
Emsleyan or Mertensian mimicry describes the unusual case where a deadly prey mimics a less dangerous species. It was first proposed by M. G. Emsley in 1966 as a possible explanation for how a predator can learn to avoid a very dangerous aposematic animal, such as a coral snake, when the predator is very likely to die, making learning unlikely. The theory was developed by the German biologist Wolfgang Wickler who named it after the German herpetologist Robert Mertens. The scenario is unlike Müllerian mimicry, where the most harmful species is the model. But if a predator dies on its first encounter with a deadly snake, it has no occasion to learn to recognize the snake's warning signals. There would then be no advantage for an extremely deadly snake in being aposematic: any predator that attacked it would be killed before it could learn to avoid the deadly prey, so the snake would be better off being camouflaged to avoid attacks. But if the predator first learnt to avoid a less deadly warning-coloured snake, the deadly species could profit by mimicking the less dangerous snake. Some harmless milk snakes (Lampropeltis triangulum), the moderately toxic false coral snakes (Erythrolamprus aesculapii), and the deadly coral snakes (Micrurus) all have a red background color with black and white/yellow rings. In this system, both the milk snakes and the deadly coral snakes are mimics, while the false coral snakes are the model.
Wasmannian
In Wasmannian mimicry, the mimic resembles a model that it lives along with in a nest or colony. Most of the models here are eusocial insects, principally ants.
Gilbertian
Gilbertian mimicry is bipolar, involving only two species. The potential host (or prey) drives away its parasite (or predator) by mimicking it, the reverse of host-parasite aggressive mimicry. It was coined by Pasteur as a phrase for such rare mimicry systems, and is named after the American ecologist Lawrence E. Gilbert who described it in 1975. The classical instance of Gilbertian mimicry is in the plant genus Passiflora, which is grazed by the micropredator larvae of some Heliconius butterflies. The host plants have evolved stipules that mimic mature Heliconius eggs near the point of hatching. The butterflies avoid laying eggs near existing ones, reducing intraspecific competition between caterpillars, which are also cannibalistic, so those that lay on vacant leaves provide their offspring with a greater chance of survival. The stipules thus appear to have evolved as Gilbertian mimics of butterfly eggs, under selection pressure from these caterpillars.
Browerian
Browerian mimicry, named after Lincoln P. Brower and Jane Van Zandt Brower who first described it in 1967, is a postulated form of automimicry; where the model belongs to the same species as the mimic. This is the analogue of Batesian mimicry within a single species, and occurs when there is a palatability spectrum within a population. Examples include the monarch and the queen from the subfamily Danainae, which feed on milkweed species of varying toxicity. These species store toxins from its host plant, which are maintained even in the adult. As levels of toxin vary depending on diet, some individuals are more toxic than the rest, which profit from the toxicity of those individuals, just as hoverflies benefit from mimicking well-defended wasps.
Misdirection by automimicry
One form of automimicry is where one part of an organism's body resembles another part. For example, the tails of some snakes resemble their heads; they move backwards when threatened and present the predator with the tail, improving their chances of escape without fatal harm. Some fishes have eyespots near their tails, and when mildly alarmed swim slowly backwards, presenting the tail as a head. Some insects such as some lycaenid butterflies have tail patterns and appendages of various degrees of sophistication that promote attacks at the rear rather than at the head. Several species of pygmy owl bear "false eyes" on the back of the head, misleading predators into reacting as though they were the subject of an aggressive stare. Many insects have filamentous "tails" at the ends of their wings and patterns of markings on the wings themselves. These combine to create a "false head". This misdirects predators such as birds and jumping spiders. Spectacular examples occur in the hairstreak butterflies; when perching on a twig or flower, they commonly do so upside down and shift their rear wings repeatedly, causing antenna-like movements of the "tails" on their wings. Studies of rear-wing damage support the hypothesis that this strategy is effective in deflecting attacks from the insect's head.
Aggressive
Predators
Aggressive mimicry is found in predators or parasites that share some of the characteristics of a harmless species, allowing them to avoid detection by their prey or host; the strategy resembles a wolf in sheep's clothing, though no conscious deceptive intent is involved. The mimic may resemble the prey or host itself, or another organism that does not threaten the prey or host.
Several spiders use aggressive mimicry to lure prey. Species such as the silver argiope (Argiope argentata) employ prominent patterns in the middle of their webs, such as zigzags. These may reflect ultraviolet light, and mimic the pattern seen in many flowers known as nectar guides. Spiders change their web day to day, which can be explained by the ability of bees to remember web patterns.
Another case is where males are lured towards what seems to be a sexually receptive female. The model in this situation is the same species as the dupe. Female fireflies of the genus Photuris emit light signals that mimic the mating signals of females of the genus Photinus. Male fireflies from several different genera are attracted to these "femmes fatales", and are captured and eaten. Each female has a repertoire of signals matching the delay and duration of the flashes of the female of the corresponding species.
Some carnivorous plants may be able to increase their rate of capturing insect prey through mimicry.
A different aggressive strategy is to mimic a mutualistic symbiont of the prey. Cleaner fish eat parasites and dead skin from client fish. Some allow the cleaner to venture inside their body to hunt these parasites. However, the sabre-toothed blenny or false cleanerfish (Aspidontus taeniatus) mimics the bluestreak cleaner wrasse (Labroides dimidiatus), which is recognized by other fishes as a cleaner. The false cleanerfish resembles the cleaner, and mimics the cleaner's "dance". Once it is allowed close to the client, it attacks, biting off a piece of its fin before fleeing. Fish wounded in this fashion soon learn to distinguish mimic from model, but because the similarity is close they also become much more cautious of the model.
A mechanism that does not involve any luring is seen in the zone-tailed hawk, which resembles the turkey vulture. It flies amongst the vultures, effectively camouflaged as a vulture which poses no threat to the hawk's prey. It hunts by suddenly breaking from the formation and ambushing its prey.
Parasites
Parasites can be aggressive mimics, though the situation is somewhat different from those outlined previously. They can mimic their hosts' natural prey, allowing themselves to be eaten as a pathway into their host. Leucochloridium, a genus of flatworm, matures in the digestive system of songbirds, their eggs then passing out of the bird in the faeces. They are then taken up by Succinea, a terrestrial snail. The eggs develop in this intermediate host, and must then find a suitable bird to mature in. Since the host birds do not eat snails, the sporocyst has another strategy to reach its host's intestine. They are brightly coloured and move in a pulsating fashion. A sporocyst-sac pulsates in the snail's eye stalks, coming to resemble an irresistible meal for a songbird. In this way, it can bridge the gap between hosts, allowing it to complete its life cycle. A nematode (Myrmeconema neotropicum) changes the colour of the abdomen of workers of the canopy ant Cephalotes atratus to make it appear like the ripe fruits of Hyeronima alchorneoides. It also changes the behaviour of the ant so that the gaster (rear part) is held raised. This presumably increases the chances of the ant being eaten by birds.
Reproductive
Reproductive mimicry occurs when the actions of the dupe directly aid in the mimic's reproduction. This is common in plants with deceptive flowers that do not provide the reward they seem to offer and it may occur in Papua New Guinea fireflies, in which the signal of Pteroptyx effulgens is used by P. tarsalis to form aggregations to attract females. Other forms of mimicry have a reproductive component, such as Vavilovian mimicry involving seeds, vocal mimicry in birds, and aggressive and Batesian mimicry in brood parasite-host systems.
Bakerian and Dodsonian
Bakerian mimicry, named after Herbert G. Baker, is a form of automimicry where female flowers mimic male flowers of their own species, cheating pollinators out of a reward. This reproductive mimicry may not be readily apparent as members of the same species may still exhibit some degree of sexual dimorphism. It is common in many species of Caricaceae.
In Dodsonian mimicry, named after Calaway H. Dodson, the model belongs to a different species than the mimic. By resembling the model, a flower can lure its pollinators without offering nectar. The mechanism occurs in several orchids, including Epidendrum ibaguense which mimics flowers of Lantana camara and Asclepias curassavica, and is pollinated by monarch butterflies and perhaps hummingbirds.
Kirbyan mimicry, brood parasitism
Brood parasitism or Kirbyan mimicry is a two species system where a brood parasite mimics its host. Cuckoos are a canonical example; the female cuckoo has its offspring raised by a bird of a different species, cutting down the biological mother's parental investment. The ability to lay eggs that mimic the host eggs is the key adaptation. The adaptation to different hosts is inherited through the female line in so-called gentes (gens, singular). Intraspecific brood parasitism, where a female lays in a conspecific's nest, as illustrated by the goldeneye duck (Bucephala clangula), do not involve mimicry The parasitic butterfly Phengaris rebeli parasitizes the ant species Myrmica schencki by releasing chemicals that fool the worker ants to believe that the caterpillar larvae are ant larvae. This enables the larvae to be brought directly into the ant's nest.
Pouyannian
In Pouyannian mimicry, a flower mimics a female of a certain insect species, inducing the males of that species to try to copulate with the flower. This is much like aggressive mimicry in fireflies, but with a more benign outcome for the pollinator. The mechanism is named after Maurice-Alexandre Pouyanne, who first described the phenomenon. It is most common in orchids, which mimic females of the order Hymenoptera (generally bees and wasps), and may account for around 60% of pollinations. Depending on the morphology of the flower, a pollen sac called a pollinium is attached to the head or abdomen of the male. This is then transferred to the stigma of the next flower the male tries to inseminate, resulting in pollination. The mimicry is a combination of visual, by olfaction, and by touch.
Vavilovian
Vavilovian mimicry is found in weeds that come to share characteristics with a domesticated plant through unintentional selection. It is named after Russian botanist and geneticist Nikolai Vavilov. Selection against the weed may occur either by manually killing the weed, or by separating its seeds from those of the crop by winnowing. Vavilovian mimicry illustrates unintentional selection by man. Weeders do not want to select weeds and their seeds that look increasingly like cultivated plants, yet there is no other option. For example, early barnyard grass, Echinochloa oryzoides, is a weed in rice fields and looks similar to rice; its seeds are often mixed in rice and have become difficult to separate through Vavilovian mimicry. Vavilovian mimics may eventually be domesticated themselves, as in the case of rye in wheat; Vavilov called these weed-crops secondary crops.
Inter-sexual mimicry
Inter-sexual mimicry (a type of automimicry, as it is within a single species) occurs when individuals of one sex in a species mimic members of the opposite sex to facilitate sneak mating. An example is the three male forms of the marine isopod Paracerceis sculpta. Alpha males are the largest and guard a harem of females. Beta males mimic females and manage to enter the harem of females without being detected by the alpha males allowing them to mate. Gamma males are the smallest males and mimic juveniles. This also allows them to mate with the females without the alpha males detecting them. Similarly, among common side-blotched lizards, some males mimic the yellow throat coloration and even mating rejection behaviour of the other sex to sneak matings with guarded females. These males look and behave like unreceptive females. This strategy is effective against "usurper" males with orange throats, but ineffective against blue throated "guarder" males, which chase them away. Female spotted hyenas have pseudo-penises that make them look like males.
| Biology and health sciences | Basics_2 | Biology |
247216 | https://en.wikipedia.org/wiki/Amaranthaceae | Amaranthaceae | Amaranthaceae ( ) is a family of flowering plants commonly known as the amaranth family, in reference to its type genus Amaranthus. It includes the former goosefoot family Chenopodiaceae and contains about 165 genera and 2,040 species, making it the most species-rich lineage within its parent order, Caryophyllales.
Description
Most species in the Amaranthaceae are annual or perennial herbs or subshrubs; others are shrubs; very few species are vines or trees. Some species are succulent. Many species have stems with thickened nodes. The wood of the perennial stem has a typical "anomalous" secondary growth; only in subfamily Polycnemoideae is secondary growth normal.
The leaves are simple and mostly alternate, sometimes opposite. They never possess stipules. They are flat or terete, and their shape is extremely variable, with entire or toothed margins. In some species, the leaves are reduced to minute scales. In most cases, neither basal nor terminal aggregations of leaves occur.
The flowers are solitary or aggregated in cymes, spikes, or panicles and typically perfect (bisexual) and actinomorphic. Some species have unisexual flowers. Bracts and bracteoles are either herbaceous or scarious. Flowers are regular with an herbaceous or scarious perianth of (one to) mostly five (rarely to eight) tepals, often joined. One to five stamens are opposite to tepals or alternating, inserting from a hypogynous disc, which may have appendages (pseudostaminodes) in some species. The anthers have two or four pollen sacs (locules). In tribe Caroxyloneae, anthers have vesicular appendages. The pollen grains are spherical with many pores (pantoporate), with pore numbers from a few to 250 (in Froelichia). One to three (rarely six) carpels are fused to a superior ovary with one (rarely two) basal ovule. Idioblasts are found in the tissues.
The diaspores are seeds or fruits (utricles), more often the perianth persists and is modified in fruit for means of dispersal. Sometimes even bracts and bracteoles may belong to the diaspore. More rarely the fruit is a circumscissile capsule or a berry. The horizontal or vertical seed often has a thickened or woody seed coat. The green or white embryo is either spirally (and without perisperm) or annular (rarely straight).
Chromosome number
The basic chromosome number is (rarely 6) mostly 8–9 (rarely 17).
Phytochemistry
Widespread in the Amaranthaceae is the occurrence of betalain pigments. The former Chenopodiaceae often contain isoflavonoids.
In phytochemical research, several methylenedioxyflavonols, saponins, triterpenoids, ecdysteroids, and specific root-located carbohydrates have been found in these plants.
Photosynthesis pathway
Although most of the family use the more common photosynthesis pathway, around 800 species are plants; this makes the Amaranthaceae the largest group with this photosynthesis pathway among the eudicots (which collectively includes about 1,600 species). Within the family, several types of photosynthesis occur, and about 17 different types of leaf anatomy are realized. Therefore, this photosynthesis pathway seems to have developed about 15 times independently during the evolution of the family. About two-thirds of the species belong to the former Chenopodiaceae. The first occurrence of photosynthesis dates from the early Miocene, about 24 million years ago, but in some groups, this pathway evolved much later, about 6 (or less) million years ago.
The multiple origin of photosynthesis in the Amaranthaceae is regarded as an evolutionary response to inexorably decreasing atmospheric levels, coupled with a more recent permanent shortage in water supply as well as high temperatures. Species that use water more efficiently had a selective advantage and were able to spread out into arid habitats.
Taxonomy
In the APG IV system of 2016, as in the previous Angiosperm Phylogeny Group classifications, the family is placed in the order Caryophyllales and includes the plants formerly treated as the family Chenopodiaceae. The monophyly of this broadly defined Amaranthaceae has been strongly supported by both morphological and phylogenetic analyses.
The family Amaranthaceae was first published in 1789 by Antoine Laurent de Jussieu in Genera Plantarum, p. 87–88. The first publication of family Chenopodiaceae was in 1799 by Étienne Pierre Ventenat in Tableau du Regne Vegetal, 2, p. 253. The older name has priority and is now the valid scientific name of the extended Amaranthaceae (s.l. = sensu lato).
Some publications still continued to use the family name Chenopodiaceae. Phylogenetic research revealed the important impact of the subfamily Polycnemoideae on the classification (see cladogram): if Polycnemoideae are considered a part of Chenopodiaceae, then Amaranthaceae (s.str. = sensu stricto) have to be included, too, and the name of the extended family is Amaranthaceae. If Polycnemoideae would be separated as its own family, Chenopodiaceae and Amaranthaceae (s.str.) would form two distinct monophyletic groups and could be treated as two separate families.
Amaranthaceae (s.l.) includes the former families Achyranthaceae , Atriplicaceae , Betaceae , Blitaceae , Celosiaceae , Chenopodiaceae nom. cons., Corispermaceae , Deeringiaceae , Dysphaniaceae nom. cons., Gomphrenaceae , Polycnemaceae , Salicorniaceae , Salsolaceae , and Spinaciaceae .
The systematics of Amaranthaceae are the subject of intensive recent research. Molecular genetic studies revealed the traditional classification, based on morphological and anatomical characters, often did not reflect the phylogenetic relationships.
The former Amaranthaceae (in their narrow circumscription) are classified into two subfamilies, Amaranthoideae and Gomphrenoideae, and contain about 65 genera and 900 species in tropical Africa and North America. The Amaranthoideae and some genera of Gomphrenoideae were found to be polyphyletic, so taxonomic changes are needed.
Current studies classified the species of former Chenopodiaceae to eight distinct subfamilies (the research is not yet completed): Polycnemoideae, which are regarded as a basal lineage, Betoideae, Camphorosmoideae, Chenopodioideae, Corispermoideae, Salicornioideae, Salsoloideae, and Suaedoideae. In this preliminary classification, the Amaranthaceae s.l. are divided into 10 subfamilies with approximately 180 genera and 2,500 species.
Genera
183 genera are accepted. A short synoptic list of genera is given here. For further and more detailed information, see the subfamily pages.
Distribution and habitat
Amaranthaceae is a widespread and cosmopolitan family from the tropics to cool temperate regions. The Amaranthaceae (sensu stricto) are predominantly tropical, whereas the former Chenopodiaceae have their centers of diversity in dry temperate and warm temperate areas. Many of the species are halophytes, tolerating salty soils, or grow in dry steppes or semi-deserts.
Uses
Some species, such as spinach (Spinacia oleracea) or forms of beet (Beta vulgaris) (beetroot, chard), are used as vegetables. Forms of Beta vulgaris include fodder beet (Mangelwurzel) and sugar beet. The seeds of Amaranthus, lamb's quarters (Chenopodium berlandieri), quinoa (Chenopodium quinoa) and kañiwa (Chenopodium pallidicaule) are edible and are used as pseudocereals.
Dysphania ambrosioides (epazote) and Dysphania anthelmintica are used as medicinal herbs. Several amaranth species are also used indirectly as a source of soda ash, such as members of the genus Salicornia (see glasswort).
A number of species are popular garden ornamental plants, especially species from the genera Alternanthera, Amaranthus, Celosia, and Iresine. Other species are considered weeds, e.g., redroot pigweed (Amaranthus retroflexus) and alligatorweed (Alternanthera philoxeroides), and several are problematic invasive species, particularly in North America, including Kali tragus and Bassia scoparia. Many species are known to cause pollen allergies.
| Biology and health sciences | Caryophyllales | Plants |
247323 | https://en.wikipedia.org/wiki/Stirling%20cycle | Stirling cycle | The Stirling cycle is a thermodynamic cycle that describes the general class of Stirling devices. This includes the original Stirling engine that was invented, developed and patented in 1816 by Robert Stirling with help from his brother, an engineer.
The ideal Otto and Diesel cycles are not totally reversible because they involve heat transfer through a finite temperature difference during the irreversible isochoric/isobaric heat-addition and heat-rejection processes. The irreversibility renders the thermal efficiency of these cycles less than that of a Carnot engine operating within the same limits of temperature. Another cycle that features isothermal heat-addition and heat-rejection processes is the Stirling cycle, which is an altered version of the Carnot cycle in which the two isentropic processes featured in the Carnot cycle are replaced by two constant-volume regeneration processes.
The cycle is reversible, meaning that if supplied with mechanical power, it can function as a heat pump for heating or cooling, and even for cryogenic cooling. The cycle is defined as a closed regenerative cycle with a gaseous working fluid. "Closed cycle" means the working fluid is permanently contained within the thermodynamic system. This also categorizes the engine device as an external heat engine. "Regenerative" refers to the use of an internal heat exchanger called a regenerator which increases the device's thermal efficiency.
The cycle is the same as most other heat cycles in that there are four main processes: compression, heat addition, expansion, and heat removal. However, these processes are not discrete, but rather the transitions overlap.
The Stirling cycle is a highly advanced subject that has defied analysis by many experts for over 190 years. Highly advanced thermodynamics is required to describe the cycle. Professor Israel Urieli writes: "...the various 'ideal' cycles (such as the Schmidt cycle) are neither physically realizable nor representative of the Stirling cycle".
The analytical problem of the regenerator (the central heat exchanger in the Stirling cycle) is judged by Jakob to rank "among the most difficult and involved that are encountered in engineering".
Idealized Stirling cycle thermodynamics
The idealized Stirling cycle consists of four thermodynamic processes acting on the working fluid (See diagram to right):
Piston motion variations
Most thermodynamics textbooks describe a highly simplified form of Stirling cycle consisting of four processes. This is known as an "ideal Stirling cycle", because it is an "idealized" model, and not necessarily an optimized cycle. Theoretically, the "ideal cycle" does have high net work output, but it is rarely used in practical applications, in part because other cycles are simpler or reduce peak stresses on bearings and other components. For convenience, the designer may elect to use piston motions dictated by system dynamics, such as mechanical linkage mechanisms. At any rate, the efficiency and cycle power are nearly as good as an actual implementation of the idealized case. A typical piston crank or linkage in a so named "kinematic" design often results in a near-sinusoidal piston motion. Some designs will cause the piston to "dwell" at either extreme of travel.
Many kinematic linkages, such as the well known "Ross yoke", will exhibit near-sinusoidal motion. However, other linkages, such as the "rhombic drive", will exhibit more non-sinusoidal motion. To a lesser extent, the ideal cycle introduces complications, since it would require somewhat higher piston acceleration and higher viscous pumping losses of the working fluid. The material stresses and pumping losses in an optimized engine, however, would only be intolerable when approaching the "ideal cycle" and/or at high cycle rates. Other issues include the time required for heat transfer, particularly for the isothermal processes. In an engine with a cycle approaching the "ideal cycle", the cycle rate might have to be reduced to address these issues.
In the most basic model of a free piston device, the kinematics will result in simple harmonic motion.
Volume variations
In beta and gamma engines, generally the phase angle difference between the piston motions is not the same as the phase angle of the volume variations. However, in the alpha Stirling, they are the same. The rest of the article assumes sinusoidal volume variations, as in an alpha Stirling with co-linear pistons, so named an "opposed piston" alpha device.
caveat: Among the many inaccuracies in this article, a co-linear alpha configuration is referenced, above. Such a configuration would be beta. Alternatively, it would be an alpha, that has an unacceptably inefficient linkage system.
Pressure-versus-volume graph
This type of plot is used to characterize almost all thermodynamic cycles. The result of sinusoidal volume variations is the quasi-elliptical shaped cycle shown in Figure 1. Compared to the idealized cycle, this cycle is a more realistic representation of most real Stirling engines. The four points in the graph indicate the crank angle in degrees.
The adiabatic Stirling cycle is similar to the idealized Stirling cycle; however, the four thermodynamic processes are slightly different (see graph above):
180° to 270°, pseudo-isothermal expansion. The expansion space is heated externally, and the gas undergoes near-isothermal expansion.
270° to 0°, near-constant-volume (or near-isometric or isochoric) heat removal. The gas is passed through the regenerator, thus cooling the gas, and transferring heat to the regenerator for use in the next cycle.
0° to 90°, pseudo-isothermal compression. The compression space is intercooled, so the gas undergoes near-isothermal compression.
90° to 180°, near-constant-volume (near-isometric or isochoric) heat addition. The compressed air flows back through the regenerator and picks up heat on the way to the heated expansion space.
With the exception of a Stirling thermoacoustic engine, none of the gas particles actually flow through the complete cycle. So this approach is not amenable to further analysis of the cycle. However, it provides an overview and indicates the cycle work.
Particle/mass motion
Figure 2 shows the streaklines which indicate how gas flows through a real Stirling engine. The vertical colored lines delineate the volumes of the engine. From left to right, they are: the volume swept by the expansion (power) piston, the clearance volume (which prevents the piston from contacting the hot heat exchanger), the heater, the regenerator, the cooler, the cooler clearance volume, and the compression volume swept by the compression piston.
Heat-exchanger pressure drop
Also referred to as "pumping losses", the pressure drops shown in Figure 3 are caused by viscous flow through the heat exchangers. The red line represents the heater, green is the regenerator, and blue is the cooler. To properly design the heat exchangers, multivariate optimization is required to obtain sufficient heat transfer with acceptable flow losses. The flow losses shown here are relatively low, and they are barely visible in the following image, which will show the overall pressure variations in the cycle.
Pressure versus crank angle
Figure 4 shows results from an "adiabatic simulation" with non-ideal heat exchangers. Note that the pressure drop across the regenerator is very low compared to the overall pressure variation in the cycle.
Temperature versus crank angle
Figure 5 illustrates the adiabatic properties of a real heat exchanger. The straight lines represent the temperatures of the solid portion of the heat exchanger, and the curves are the gas temperatures of the respective spaces. The gas temperature fluctuations are caused by the effects of compression and expansion in the engine, together with non-ideal heat exchangers which have a limited rate of heat transfer. When the gas temperature deviates above and below the heat exchanger temperature, it causes thermodynamic losses known as "heat transfer losses" or "hysteresis losses". However, the heat exchangers still work well enough to allow the real cycle to be effective, even if the actual thermal efficiency of the overall system is only about half of the theoretical limit.
Cumulative heat and work energy
Figure 6 shows a graph of the alpha-type Stirling engine data, where 'Q' denotes heat energy, and 'W' denotes work energy. The blue dotted line shows the work output of the compression space. As the trace dips down, work is done on the gas as it is compressed. During the expansion process of the cycle, some work is actually done on the compression piston, as reflected by the upward movement of the trace. At the end of the cycle, this value is negative, indicating that compression piston requires a net input of work. The blue solid line shows the heat flowing out of the cooler heat exchanger. The heat from the cooler and the work from the compression piston have the same cycle energy. This is consistent with the zero-net heat transfer of the regenerator (solid green line). As would be expected, the heater and the expansion space both have positive energy flow. The black dotted line shows the net work output of the cycle. On this trace, the cycle ends higher than it started, indicating that the heat engine converts energy from heat into work.
| Physical sciences | Thermodynamics | Physics |
247326 | https://en.wikipedia.org/wiki/Volcanic%20glass | Volcanic glass | Volcanic glass is the amorphous (uncrystallized) product of rapidly cooling magma. Like all types of glass, it is a state of matter intermediate between the closely packed, highly ordered array of a crystal and the highly disordered array of liquid. Volcanic glass may refer to the interstitial material, or matrix, in an aphanitic (fine-grained) volcanic rock, or to any of several types of vitreous igneous rocks.
Origin
Volcanic glass is formed when magma is rapidly cooled. Magma rapidly cooled to below its normal crystallization temperature becomes a supercooled liquid, and, with further rapid cooling, this becomes an amorphous solid. The change from supercooled liquid to glass occurs at a temperature called the glass transition temperature, which depends on both cooling rate and the amount of water dissolved in the magma. Magma rich in silica and poor in dissolved water is most easily cooled rapidly enough to form volcanic glass. As a result, rhyolite magmas, which are high in silica, can produce tephra composed entirely of volcanic glass and may also form glassy lava flows. Ash-flow tuffs typically consist of countless microscopic shards of volcanic glass. Basalt, which is low in silica, forms glass only with difficulty, so that basalt tephra almost always contains at least some crystalline material (quench crystals). The glass transition temperature of basalt is about .
The mechanisms controlling formation of volcanic glass are further illustrated by the two forms of basaltic glass, tachylite and sideromelane. Tachylite is opaque to transmitted light because of the abundance of tiny oxide mineral crystals suspended in the glass. Sideromelane is partially transparent because it contains much fewer crystals. Sideromelane is abundant only in eruptions where basalt magma has been very rapidly cooled by contact with water, such as phreatomagmatic eruptions. Basaltic volcanic glass is also present in pillow lavas.
Of the cooling mechanisms responsible for forming volcanic glass, the most effective is quenching by water, followed by cooling by entrained air in an eruption column. The least effective mechanism is cooling at the bottom of a flow in contact with the ground.
Types
Most commonly, volcanic glass refers to obsidian, a rhyolitic glass with high silica (SiO2) content.
Other types of volcanic glass include the following:
Pumice, which is considered a glass because it has no crystal structure.
Apache tears, a kind of nodular obsidian.
Tachylite (also spelled tachylyte), a basaltic glass with relatively low silica content.
Sideromelane, a less common form of tachylyte.
Palagonite, an alteration product of basaltic glass.
Hyaloclastite, a hydrated tuff-like breccia of sideromelane and palagonite.
Pele's hair, threads or fibers of volcanic glass, usually basaltic.
Pele's tears, tear-like drops of volcanic glass, usually basaltic.
Limu o Pele (Pele's seaweed), thin sheets and flakes of brownish-green to near-clear volcanic glass, usually basaltic.
Alteration
Volcanic glass is chemically unstable and readily decomposes. Water molecules readily react with the open, disordered structure of volcanic glass, removing soluble cations from the glass and precipitating secondary (authigenic) minerals. As a result, lithification of volcanic ash is one of the fastest low-temperature lithification processes. Alteration of volcanic glass at mid-ocean ridges may have contributed significantly to the formation of massive sulfide deposits, and alteration of volcanic ash beds formed economically important zeolite and bentonite deposits.
| Physical sciences | Igneous rocks | Earth science |
247388 | https://en.wikipedia.org/wiki/Combine%20harvester | Combine harvester | The modern combine harvester, also called a combine, is a machine designed to harvest a variety of cultivated seeds. Combine harvesters are one of the most economically important labour-saving inventions, significantly reducing the fraction of the population engaged in agriculture. Among the crops harvested with a combine are wheat, rice, oats, rye, barley, corn (maize), sorghum, millet, soybeans, flax (linseed), sunflowers and rapeseed (canola). The separated straw (consisting of stems and any remaining leaves with limited nutrients left in it) is then either chopped onto the field and ploughed back in, or laid out in rows, ready to be baled and used for bedding and cattle feed.
The name of the machine is derived from the fact that the harvester combined multiple separate harvesting operations – reaping, threshing or winnowing and gathering – into a single process around the start of the 20th century. A combine harvester still performs those operation principles. The machine can easily be divided into four parts, namely: the intake mechanism, the threshing and separation system, the cleaning system, and finally the grain handling and storage system. Electronic monitoring assists the operator by providing an overview of the machine's operation, and the field's yield.
History
In 1826 in Scotland, the inventor Reverend Patrick Bell designed a reaper machine, which used the scissors principle of plant cutting (a principle that is used to this day). The Bell machine was pushed by horses. A few Bell machines were available in the United States. In 1835, in the United States, Hiram Moore built and patented the first combine harvester, which was capable of reaping, threshing and winnowing cereal grain. Early versions were pulled by horse, mule or ox teams. In 1835, Moore built a full-scale version with a length of 5.2 m (17 ft) and a cut width of 4.57 m (15 ft); by 1839, over of crops were harvested. This combine harvester was pulled by 20 horses fully handled by farmhands. By 1860, combine harvesters with a cutting, or swathe, width of several metres were used on American farms.
A parallel development in Australia saw the development of the stripper based on the Gallic stripper, by John Ridley and others in South Australia by 1843. The stripper only gathered the heads, leaving the stems in the field. The stripper and later headers had the advantage of fewer moving parts and only collecting heads, requiring less power to operate. Refinements by Hugh Victor McKay produced a commercially successful combine harvester in 1885, the Sunshine Header-Harvester.
Combines, some of them quite large, were drawn by mule or horse teams and used a bullwheel to provide power. Later, steam power was used, and George Stockton Berry integrated the combine with a steam engine using straw to heat the boiler. At the turn of the twentieth century, horse-drawn combines were starting to be used on the American plains and Idaho (often pulled by teams of twenty or more horses).
In 1911, the Holt Manufacturing Company of California, US produced a self-propelled harvester. In Australia in 1923, the patented Sunshine Auto Header was one of the first center-feeding self-propelled harvesters. In 1923 in Kansas, the Baldwin brothers and their Gleaner Manufacturing Company patented a self-propelled harvester that included several other modern improvements in grain handling. Both the Gleaner and the Sunshine used Fordson engines; early Gleaners used the entire Fordson chassis and driveline as a platform. In 1929, Alfredo Rotania of Argentina patented a self-propelled harvester.
International Harvester started making horse-pulled combines in 1915. At the time, horse-powered binders and stand-alone threshing machines were more common. In the 1920s, Case Corporation and John Deere made combines, introducing tractor-pulled harvesters with a second engine aboard the combine to power its workings. The world economic collapse in the 1930s stopped farm equipment purchases, and for this reason, people largely retained the older method of harvesting. A few farms did invest and used Caterpillar tractors to move the outfits.
Tractor-drawn combines (also called pull-type combines) became common after World War II as many farms began to use tractors. An example was the All-Crop Harvester series. These combines used a shaker to separate the grain from the chaff and straw-walkers (grates with small teeth on an eccentric shaft) to eject the straw while retaining the grain. Early tractor-drawn combines were usually powered by a separate gasoline engine, while later models were PTO-powered, via a shaft transferring tractor engine power to operate the combine. These machines either put the harvested crop into bags that were then loaded onto a wagon or truck, or had a small bin that stored the grain until it was transferred via a chute.
In the U.S., Allis-Chalmers, Massey-Harris, International Harvester, Gleaner Manufacturing Company, John Deere, and Minneapolis Moline are past or present major combine producers.
In 1937, the Australian-born Thomas Carroll, working for Massey-Harris in Canada, perfected a self-propelled model and in 1940, a lighter-weight model began to be marketed widely by the company. Lyle Yost invented an auger that would lift grain out of a combine in 1947, making unloading grain much easier and further from the combine. In 1952 Claeys launched the first self-propelled combine harvester in Europe; in 1953, the European manufacturer Claas developed a self-propelled combine harvester named 'Hercules', it could harvest up to 5 tons of wheat a day. This newer kind of combine is still in use and is powered by diesel or gasoline engines. Until the self-cleaning rotary screen was invented in the mid-1960s combine engines suffered from overheating as the chaff spewed out when harvesting small grains would clog radiators, blocking the airflow needed for cooling.
A significant advance in the design of combines was the rotary design. The grain is initially stripped from the stalk by passing along a helical rotor, instead of passing between rasp bars on the outside of a cylinder and a concave. Rotary combines were first introduced by Sperry-New Holland in 1975.
Around the 1980s, on-board electronics were introduced to measure threshing efficiency. This new instrumentation allowed operators to get better grain yields by optimizing ground speed and other operating parameters.
The largest "class 10-plus" combines, which emerged in the early 2020's, have nearly 800 engine horsepower (600 kW) and are fitted with headers up to wide.
Combine header
Combines are equipped with removable headers that are designed for particular crops. The standard header, sometimes called a grain platform, is equipped with a reciprocating knife cutter bar, and features a revolving reel with metal teeth to cause the cut crop to fall into the auger once it is cut. A variation of the platform, a "flex" platform, is similar but has a cutter bar that can flex over contours and ridges to cut soybeans that have pods close to the ground. A flex head can cut soybeans as well as cereal crops, while a rigid platform is generally used only in cereal grains.
Some wheat headers, called "draper" headers, use a fabric or rubber apron instead of a cross auger. Draper headers allow faster feeding than cross augers, leading to higher throughputs due to lower power requirements. On many farms, platform headers are used to cut wheat, instead of separate wheat headers, so as to reduce overall costs.
Dummy heads or pick-up headers feature spring-tined pickups, usually attached to a heavy rubber belt. They are used for crops that have already been cut and placed in windrows or swaths. This is particularly useful in northern climates such as western Canada, where swathing kills weeds resulting in a faster dry down.
While a grain platform can be used for corn, a specialized corn head is ordinarily used instead. The corn head is equipped with snap rolls that strip the stalk and leaf away from the ear, so that only the ear (and husk) enter the throat. This improves efficiency dramatically since so much less material must go through the cylinder. The corn head can be recognized by the presence of points between each row.
Occasionally rowcrop heads are seen that function like a grain platform but have points between rows like a corn head. These are used to reduce the amount of weed seed picked up when harvesting small grains.
Self-propelled Gleaner combines could be fitted with special tracks instead of tires to assist in harvesting rice. These tracks can be made to fit other combines by adding adapter plates. Some combines, particularly the pull type, have tires with a deep diamond tread which prevents sinking in mud.
Conventional combine
The cut crop is carried up the feeder throat (commonly called the "feederhouse"), by a chain and flight elevator, then fed into the threshing mechanism of the combine, consisting of a rotating threshing drum (commonly called the "cylinder"), to which grooved steel bars (rasp bars) are bolted. The rasp bars thresh or separate the grains and chaff from the straw through the action of the cylinder against the concave, a shaped "half drum", also fitted with steel bars and a meshed grill, through which grain, chaff and smaller debris may fall, whereas the straw, being too long, is carried through onto the straw walkers. This action is also allowed because grain is heavier than straw, which causes it to fall rather than "float" across from the cylinder/concave to the walkers. The drum speed is variably adjustable on most machines, whilst the distance between the drum and concave is finely adjustable fore, aft and together, to achieve optimum separation and output. Manually engaged disawning plates are usually fitted to the concave. These provide extra friction to remove the awns from barley crops. After the primary separation at the cylinder, the clean grain falls through the concave and to the shoe, which contains the chaffer and sieves. The shoe is common to both conventional combines and rotary combines.
Hillside leveling
Hillside leveling, in which a hydraulic system re-orients the combine, allows combines to harvest steep but fertile soil. Their primary advantage is increased threshing efficiency. Without leveling, grain and chaff slide to one side of separator and come through the machine in a large ball rather than being separated, dumping large amounts of grain on the ground. By keeping the machinery level, the straw-walker is able to thresh more efficiently. Secondarily, leveling changes a combine's center of gravity relative to the hill and allows the combine to harvest along the contour of a hill without tipping, a danger on steeper slopes; it is not uncommon for combines to roll over on extremely steep hills. Hillside leveling can be very important in regions with steep hills, such as the Palouse region of the Pacific Northwest of the United States, where hillsides can have slopes as steep as 50%.
The first leveling technology was developed by Holt Co., a US company in California, in 1891. Modern leveling came into being with the invention and patent of a level sensitive mercury switch system invented by Raymond Alvah Hanson in 1946. A leveling system was also developed in Europe by the Italian combine manufacturer Laverda. Gleaner, IH/Case IH, John Deere, and others all have made combines with a hillside leveling system, and local machine shops have fabricated them as an aftermarket add-on. Newer leveling systems do not have as much tilt as the older ones, as modern combines use a rotary grain separator which makes leveling less critical.
Sidehill leveling
Sidehill combines are very similar to hillside combines in that they level the combine to the ground so that the threshing can be efficiently conducted; however, they have some very distinct differences. Modern hillside combines level around 35% on average, while older machines were closer to 50%. Sidehill combines only level to 18%. They are sparsely used in the Palouse region. Rather, they are used on the gentle rolling slopes of the midwest. Sidehill combines are much more mass-produced than their hillside counterparts. The height of a sidehill machine is the same height as a level-land combine. Hillside combines have added steel that sets them up approximately 2–5 feet higher than a level-land combine and provide a smooth ride.
Maintaining threshing speed
Another technology that is sometimes used on combines is a continuously variable transmission. This allows the ground speed of the machine to be varied while maintaining a constant engine and threshing speed. It is desirable to keep the threshing speed constant since the machine will typically have been adjusted to operate best at a certain speed.
Self-propelled combines started with standard manual transmissions that provided one speed based on input rpm. Deficiencies were noted and in the early 1950s combines were equipped with what John Deere called the "Variable Speed Drive". This was simply a variable width sheave controlled by spring and hydraulic pressures. This sheave was attached to the input shaft of the transmission. A standard 4-speed manual transmission was still used in this drive system. The operator would select a gear, typically 3rd. An extra control was provided to the operator to allow him to speed up and slow down the machine within the limits provided by the variable speed drive system. By decreasing the width of the sheave on the input shaft of the transmission, the belt would ride higher in the groove. This slowed the rotating speed on the input shaft of the transmission, thus slowing the ground speed for that gear. A clutch was still provided to allow the operator to stop the machine and change transmission gears.
Later, as hydraulic technology improved, hydrostatic transmissions were introduced for use on swathers but later this technology was applied to combines as well. This drive retained the 4-speed manual transmission as before, but used a system of hydraulic pumps and motors to drive the input shaft of the transmission. The engine turns the hydraulic pump capable of pressures up to . This pressure is then directed to the hydraulic motor that is connected to the input shaft of the transmission. The operator is provided with a lever in the cab that allows for the control of the hydraulic motor's ability to use the energy provided by the pump.
Most if not all modern combines are equipped with hydrostatic drives. These are larger versions of the same system used in consumer and commercial lawn mowers that most are familiar with today. In fact, it was the downsizing of the combine drive system that placed these drive systems into mowers and other machines.
The operating principle
Despite great advances in mechanics and computer control, the basic operation of the combine harvester has remained unchanged almost since it was invented.
Power requirements over the years have increased due to larger capacities and some processes such as rotary threshing and straw chopping take considerable power. This is sometimes supplied by a large tractor in a pull-type combine, or a large gasoline or diesel engine in a self-propelled type. A frequent problem is the presence of airborne chaff and straw, which can accumulate causing a fire hazard and to radiators which can become plugged. Most machines have addressed these problems with enclosed engine compartments and rotary centrifugal inlet screens which prevent chaff buildup.
First, the header, described above, cuts the crop and feeds it into the threshing cylinder. This consists of a series of horizontal rasp bars fixed across the path of the crop and in the shape of a quarter cylinder. Moving rasp bars or rub bars pull the crop through concaved grates that separate the grain and chaff from the straw. The grain heads fall through the fixed concaves. What happens next is dependent on the type of combine in question. In most modern combines, the grain is transported to the shoe by a set of 2, 3, or 4 (possibly more on the largest machines) augers, set parallel or semi-parallel to the rotor on axial mounted rotors and perpendicular on axial-flow combines.
In older Gleaner machines, these augers were not present. Those combines are unique in that the cylinder and concave is set inside feederhouse instead of in the machine directly behind the feederhouse. Consequently, the material was moved by a "raddle chain" from underneath the concave to the walkers. The clean grain fell between the raddle and the walkers onto the shoe, while the straw, being longer and lighter, floated across onto the walkers to be expelled. On most other older machines, the cylinder was placed higher and farther back in the machine, and the grain moved to the shoe by falling down a "clean grain pan", and the straw "floated" across the concaves to the back of the walkers.
Since the Sperry-New Holland TR70 twin-rotor combine came out in 1975, most manufacturers have combines with rotors in place of conventional cylinders. However, makers have now returned to the market with conventional models alongside their rotary line-up. A rotor is a long, longitudinally mounted rotating cylinder with plates similar to rub bars (except for in the above-mentioned Gleaner rotaries).
There are usually two sieves, one above the other. The sieves are basically metal frames that have many rows of "fingers" set reasonably close together. The angle of the fingers is adjustable, to change the clearance and thereby control the size of material passing through. The top is set with more clearance than the bottom to allow a gradual cleaning action. Setting the concave clearance, fan speed, and sieve size is critical to ensure that the crop is threshed properly, the grain is clean of debris, and all of the grain entering the machine reaches the grain tank or 'hopper'. (Observe, for example, that when travelling uphill the fan speed must be reduced to account for the shallower gradient of the sieves.)
Heavy material, e.g., unthreshed heads, fall off the front of the sieves and are returned to the concave for re-threshing.
The straw walkers are located above the sieves, and also have holes in them. Any grain remaining attached to the straw is shaken off and falls onto the top sieve.
When the straw reaches the end of the walkers it falls out the rear of the combine. It can then be baled for cattle bedding or spread by two rotating straw spreaders with rubber arms. Most modern combines are equipped with a straw spreader.
Rather than immediately falling out the rear of the combine at the end of the walkers, there are models of combine harvesters from Eastern Europe and Russia (e.g. Agromash Yenisei 1200 1 HM, etc.) that have "straw catchers" at the end of the walkers, which temporarily hold the straw and then, once full, deposit it in a stack for easy gathering.
Thresher designs
For some time, combine harvesters used the conventional design, which used a rotating cylinder at the front-end which knocked the seeds out of the heads, and then used the rest of the machine to separate the straw from the chaff, and the chaff from the grain. The TR70 from Sperry-New Holland was brought out in 1975 as the first rotary combine. Other manufacturers soon followed, International Harvester with their "Axial-Flow" in 1977 and Gleaner with their N6 in 1979.
In the decades before the widespread adoption of the rotary combine in the late seventies, several inventors had pioneered designs which relied more on centrifugal force for grain separation and less on gravity alone. By the early eighties, most major manufacturers had settled on a "walkerless" design with much larger threshing cylinders to do most of the work. Advantages were faster grain harvesting and gentler treatment of fragile seeds, which were often cracked by the faster rotational speeds of conventional combine threshing cylinders.
It was the disadvantages of the rotary combine (increased power requirements and over-pulverization of the straw by-product) which prompted a resurgence of conventional combines in the late nineties. Perhaps overlooked but nonetheless true, when the large engines that powered the rotary machines were employed in conventional machines, the two types of machines delivered similar production capacities. Also, research was beginning to show that incorporating above-ground crop residue (straw) into the soil is less useful for rebuilding soil fertility than previously believed. This meant that working pulverized straw into the soil became more of a hindrance than a benefit. An increase in feedlot beef production also created a higher demand for straw as fodder. Conventional combines, which use straw walkers, preserve the quality of straw and allow it to be baled and removed from the field.
Instrumentation
While the principles of basic threshing have changed little over the years, modern advancements in electronics and monitoring technology has continued to develop. Whereas older machines required the operator to rely on machine knowledge, frequent inspection and monitoring, and a keen ear to listen for subtle sound changes, newer machines have replaced many of those duties with instrumentation.
Shaft monitors
Early on, simple magnetic pickups were used to monitor shaft rotation, and issue a warning when they deviated beyond preset limits. Temperature sensors can also give warning when bearings overheat due to lack of lubrication, sometimes leading to combine fires.
Loss monitors
The job of monitoring how much grain is wasted by the thresher by being discharged with the chaff and straw used to require going behind the machine to check. Yield monitors work like a microphone, registering an electrical impulse caused by grains impacting a plate. A meter in the operator's cab displays the relative amount of grain loss proportional to speed.
Yield monitoring
Measuring the amount of yield (bushels per acre or tonnes per hectare) has become increasingly important, particularly when real-time measurement can help determine which areas of a field are more or less productive. These variations can often be remediated with variable crop inputs. Yield is determined by measuring the amount of grain harvested in relation to the area covered.
Cameras
Cameras placed at strategic points on the machine can eliminate some of the guesswork for the operator.
Field mapping
The advent of GPS and GIS technologies has made it possible to create field maps, which can assist in navigation, and in the preparation of yield maps, which show which parts of the field are more productive.
Combine size classification
While all combines aim to achieve the same result, each machine can be classified based on its general throughput which is based upon the rated horsepower rating of the combine. Currently combine classifications, as defined by Association of Equipment Manufacturers (AEM), are as follows (metric horsepower, which is approximately 735.5 watts, is used)
Class 5 - less than 280 HP
Class 6 - 280 HP - 360 HP
Class 7 - 360 HP - 500 HP
Class 8 - 500 HP - 600 HP
Class 9 - 600 HP - 680 HP
Class 10 - more than 680 HP
While this classification is current, the classes themselves have and will evolve over time. For instance, a class 7 combine in the year 1980 would only have 270 horsepower and been one of the largest machines available in the world at that time but in the 21st century the same machine would be considered small. The Association of Equipment Manufacturers recognizes Class 10, which came into being in 2013, as the largest combine class. However, there are combines with horsepower and threshing capacity that could argue for creating a new class.
Combine fires
Grain combine fires are responsible for millions of dollars of loss each year. Fires usually start near the engine where dust and dry crop debris accumulate. Fires can also start when heat is introduced by bearings or gearboxes that have failed. From 1984 to 2000, 695 major grain combine fires were reported to U.S. local fire departments. Dragging chains to reduce static electricity was one method employed for preventing harvester fires, but it is not yet clear what if any role static electricity plays in causing harvester fires. The application of appropriate synthetic greases will reduce the friction experienced at crucial points (i.e., chains, sprockets and gear boxes) compared to petroleum based lubricants. Engines with synthetic lubricants will also remain significantly cooler during operation.
Conversions
Obsolete or damaged combines can be converted into general utility tractors. This is possible if the relevant systems (cabin, drivetrain, controls and hydraulics) still work or can be repaired. Conversions typically involve removing specialized components for threshing and processing crops; they can also include modifying the frame and controls to better suit operation as a tractor (including lowering it closer to the ground).
Thresher drives can sometimes be repurposed as power take-offs.
| Technology | Farm and garden machinery | null |
247405 | https://en.wikipedia.org/wiki/Crossbill | Crossbill | Crossbills are birds of the genus Loxia within the finch family (Fringillidae), with six species. These birds are characterized by the mandibles with crossed tips, which gives the group its English name. Adult males tend to be red or orange in color, and females green or yellow, but there is much variation.
Crossbills are specialist feeders on conifer cones, and the unusual bill shape is an adaptation which enables them to extract seeds from cones. These birds are typically found in higher northern hemisphere latitudes, where their food sources grow. They irrupt out of the breeding range when the cone crop fails. Crossbills breed very early in the year, often in winter months, to take advantage of maximum cone supplies.
Systematics and evolution
The genus Loxia was introduced by the Swedish naturalist Carl Linnaeus in 1758 in the 10th edition of his Systema Naturae. The name is from the Ancient Greek , "crosswise". The Swiss naturalist Conrad Gessner had used the word Loxia for a crossbill in 1555 in his Historiae Animalium. The type species was designated as Loxia curvirostra (red crossbill) by George Robert Gray in 1840.
Analysis of mitochondrial cytochrome b sequence data indicates that the crossbills and redpolls share a common ancestor and only diverged during the Tortonian (, Late Miocene). The research suggests that the genera Loxia and Carduelis might be merged into a single genus, for which the name Loxia would then have priority. But this would imply changing the name of a large number of species, and given that the adaptations of the crossbills represent a unique evolutionary path (see Evolutionary grade), it seems more appropriate to split up the genus Carduelis as it had already been done during most of the 20th century. Unfortunately, the fossil record is restricted to a Late Pliocene () species, Loxia patevi, found at Varshets, Bulgaria.
The species of crossbills are difficult to separate, and care is needed even with the two-barred and Hispaniolan crossbills, the easiest. The other species are identified by subtle differences in head shape and bill size, and the identification problems formerly led to much taxonomic speculation, with some scientists considering that the parrot and Scottish crossbills and possibly the Hispaniolan and two-barred crossbills are conspecific.
The identification problem is least severe in North America, where only the red and white-winged species occur, and (possibly) worst in the Scottish Highlands, where three species breed and the two-barred is also a possible vagrant.
Work on vocalization in North America suggests that there are eight or nine discrete populations of red crossbill in that continent alone, which do not interbreed and are (like the named species) adapted to specialize in different conifer species. While several ornithologists seem inclined to give these forms species status, no division of the American red crossbills has yet occurred. Preliminary investigations in Europe and Asia suggest an equal, if not greater, complexity, with several different call types identified; these call types being as different from each other as from the named species of the parrot and Scottish crossbills - suggesting either that they are valid species, or else that the parrot and Scottish crossbills may not be.
Genetic research on their DNA failed to reveal any difference between any of the crossbills (including the morphologically distinct two-barred), with variation between individuals greater than any difference between the taxa. This led to the suggestion that limited interbreeding between the different types prevented significant genetic differentiation, and enabled each type to maintain a degree of morphological plasticity, which may be necessary to enable them to feed on different conifers when their preferred food species has a crop failure. Research in Scotland, however, has shown that the parrot and Scottish crossbills are reproductively isolated from each other and also from the red crossbill, despite irruption of that species into their ranges, and the diagnostic calls and bill dimensions have not been lost. They are, therefore, good species.
Currently accepted species and their preferred food sources are:
Originally, the chestnut-backed sparrow-lark (Eremopterix leucotis) and Pine grosbeak (Pinicola enucleator) were also classified as belonging within the genus Loxia.
Feeding behavior
The different species specialize in feeding on different conifer species, with the bill shape optimized for opening that species of conifer. This is achieved by inserting the bill between the conifer cone scales and twisting the lower mandible towards the side to which it crosses, enabling the bird to extract the seed at the bottom of the scale with its tongue.
The mechanism by which the bill-crossing is developed (which usually, but not always, occurs in a 1:1 frequency of left-crossing or right-crossing morphs), and what determines the direction, has hitherto withstood all attempts to resolve it.
It is very probable that there is a genetic basis underlying the phenomenon (young birds whose bills are still straight will give a cone-opening behavior if their bills are gently pressed, and the crossing develops before the birds are fledged and feeding independently), but at least in the red crossbill (the only species which has been somewhat thoroughly researched regarding this question) there is no straightforward mechanism of heritability.
While the direction of crossing seems to be the result of at least three genetic factors working together in a case of epistasis and most probably autosomal, it is not clear whether the 1:1 frequency of both morphs in most cases is the result of genetics or environmental selection. Populations that feed on cones without removing or twisting them will likely show a 1:1 morph distribution, no matter what the genetic basis may be: the fitness of each morph is inversely proportional to its frequency in the population. Such birds can only access the cone with the lower mandible tip pointing towards it to successfully extract seeds, and thus a too high number of birds of one morph will result in the food availability for each bird of this morph decreasing.
They can utilize other conifers to their preferred, and often need to do so when their preferred species has a crop failure, but are less efficient in their feeding (not enough to prevent survival, but probably enough to reduce breeding success).
Fossil record
Loxia patevi was described from the late Pliocene of Varshets, Bulgaria.
| Biology and health sciences | Passerida | Animals |
247410 | https://en.wikipedia.org/wiki/Cynodontia | Cynodontia | Cynodontia () is a clade of eutheriodont therapsids that first appeared in the Late Permian (approximately 260 mya), and extensively diversified after the Permian–Triassic extinction event. Mammals are cynodonts, as are their extinct ancestors and close relatives (Mammaliaformes), having evolved from advanced probainognathian cynodonts during the Late Triassic.
Non-mammalian cynodonts occupied a variety of ecological niches, both as carnivores and as herbivores. Following the emergence of mammals, most other cynodont lines went extinct, with the last known non-mammaliaform cynodont group, the Tritylodontidae, having its youngest records in the Early Cretaceous.
Description
Early cynodonts have many of the skeletal characteristics of mammals. The teeth were fully differentiated and the braincase bulged at the back of the head. Outside of some crown-group mammals (notably the therians), all cynodonts probably laid eggs. The temporal fenestrae were much larger than those of their ancestors, and the widening of the zygomatic arch in a more mammal-like skull would have allowed for more robust jaw musculature. They also have the secondary palate that other primitive therapsids lacked, except the therocephalians, who were the closest relatives of cynodonts. (However, the secondary palate of cynodonts primarily comprises the maxillae and palatines as in mammals, whereas the secondary palate of the therocephalians primarily comprises the maxillae and the vomer.) The dentary was the largest bone in their lower jaw.
The cynodonts probably had some form of warm-blooded metabolism. This has led to many reconstructions of cynodonts as having fur. Being endothermic they may have needed it for thermoregulation, but fossil evidence of their fur (or lack thereof) has been elusive. Modern mammals have Harderian glands secreting lipids to coat their fur, but the telltale imprint of this structure is only found from the primitive mammal Morganucodon and onwards. Nonetheless, recent studies on Permian synapsid coprolites show that more basal therapsids may have had fur, and at any rate fur was already present in Mammaliaformes such as Castorocauda and Megaconus.
Early cynodonts had numerous small foramina on their snout bones, similar to reptiles. This suggests that they had immobile, non-muscular lips like those of lizards, and lacked muscular cheeks. As a muscular, mobile face is necessary to perform whisking movements and to avoid damage to whiskers, it is unlikely that early cynodonts had whiskers. In prozostrodontian cynodonts, the group that includes mammals, the foramina are replaced by a single large infraorbital foramen, which indicates that the face had become muscular and that whiskers would have been present.
Derived cynodonts developed epipubic bones. These served to strengthen the torso and support abdominal and hindlimb musculature, aiding them in the development of an erect gait, but at the expense of prolonged pregnancy, forcing these animals to give birth to highly altricial young as in modern marsupials and monotremes. Only placentals, and perhaps Megazostrodon and Erythrotherium, would lose these. A specimen of Kayentatherium does indeed demonstrate that at least tritylodontids already had a fundamentally marsupial-like reproductive style, but produced much higher litters at around 38 perinates or possibly eggs.
Cynodonts are the only known synapsid lineage to have produced aerial locomotors, with gliding being known in haramiyidans and various mammal groups, and placental mammals having developed flight.
The largest known non-mammalian cynodont is Scalenodontoides, a traversodontid, which has been estimated to have a maximum skull length of approximately based on a fragmentary specimen.
Evolutionary history
The closest relatives of cynodonts are therocephalians, with which they form the clade Eutheriodontia.
The earliest cynodonts are known early Lopingian (early Wuchiapingian) aged sediments of the Tropidostoma Assemblage Zone, in the Karoo Supergroup of South Africa, belonging to the basal family Charassognathidae. Fossils of Permian cynodonts are relatively rare outside of South Africa, with the most widely distributed genus being Procynosuchus, which is known from South Africa, Germany, Tanzania, Zambia, and possibly Russia.
Cynodonts expanded rapidly in diversity after the Permian–Triassic extinction event. Peak disparity in cynodonts occurred from the Induan to the Carnian and in the middle Norian. Post-Early Triassic cynodonts were dominated by members of the advanced clade Eucynodontia, which has two main subdivisions, the predominantly herbivorous Cynognathia and the predominantly carnivorous Probainognathia. During the Early and Middle Triassic, cynodont diversity was dominated by members of Cynognathia, and members of Probainognathia would not become prominent until the Late Triassic (early Norian). Almost all Middle Triassic cynodonts are known from Gondwana, with only one genus (Nanogomphodon) having been found in the Northern Hemisphere. Among the most dominant groups of Middle and Late Triassic cynodonts is the herbivorous Traversodontidae, predominantly in Gondwana, which reached a peak diversity in the Late Triassic. Mammaliaformes originated from probainognathian cynodonts during the Late Triassic. Early Mammaliaformes were small bodied insectivores. Only two groups of non-mammaliaform cynodonts existed beyond the end of the Triassic, both belonging to Probainognathia. The first is the insectivorous Tritheledontidae, which briefly lasted into the Early Jurassic. The second is the herbivorous Tritylodontidae, which first appeared in the latest Triassic, which were abundant and diverse during the Jurassic, predominantly in the Northern Hemisphere, persisted into the Early Cretaceous (Barremian-Aptian) in Asia, at least until around 120 million years ago, as represented by Fossiomanus from China.
During their evolution, the number of cynodont jaw bones reduced. This move towards a single bone for the mandible paved the way for other bones in the jaw, the articular and angular, to migrate to the cranium, where they function as parts of the mammalian hearing system.
Cynodonts also developed a secondary palate in the roof of the mouth. This caused air flow from the nostrils to travel to a position in the back of the mouth instead of directly through it, allowing cynodonts to chew and breathe at the same time. This characteristic is present in all mammals.
Taxonomy
Richard Owen named Cynodontia in 1861, which he assigned to Anomodontia as a family. Robert Broom (1913) reranked Cynodontia as an infraorder, since retained by others, including Colbert and Kitching (1977), Carroll (1988), Gauthier et al. (1989), and Rubidge and Cristian Sidor (2001). Olson (1966) assigned Cynodontia to Theriodontia, Colbert and Kitching (1977) to Theriodontia, and Rubridge and Sidor (2001) to Eutheriodontia. William King Gregory (1910), Broom (1913), Carroll (1988), Gauthier et al. (1989), Hopson and Kitching (2001) and Botha et al. (2007) all considered Cynodontia as belonging to Therapsida. Botha et al. (2007) seems to have followed Owen (1861), but without specifying taxonomic rank.
Phylogeny
Below is a cladogram from Ruta, Botha-Brink, Mitchell and Benton (2013) showing one hypothesis of cynodont relationships:
Distribution
Non-mammalian cynodonts have been found in South America, India, Africa, Antarctica, Asia, Europe and North America.
| Biology and health sciences | Proto-mammals | Animals |
247414 | https://en.wikipedia.org/wiki/Y-chromosomal%20Adam | Y-chromosomal Adam | In human genetics, the Y-chromosomal Adam (more technically known as the Y-chromosomal most recent common ancestor, shortened to Y-MRCA), is the patrilineal most recent common ancestor (MRCA) from whom all currently living humans are descended. He is the most recent male from whom all living humans are descended through an unbroken line of their male ancestors. The term Y-MRCA reflects the fact that the Y chromosomes of all currently living human males are directly derived from the Y chromosome of this remote ancestor.
The analogous concept of the matrilineal most recent common ancestor is known as "Mitochondrial Eve" (mt-MRCA, named for the matrilineal transmission of mtDNA), the most recent woman from whom all living humans are descended matrilineally. As with "Mitochondrial Eve", the title of "Y-chromosomal Adam" is not permanently fixed to a single individual, but can advance over the course of human history as paternal lineages become extinct.
Estimates of the time when Y-MRCA lived have also shifted as modern knowledge of human ancestry changes. For example, in 2013, the discovery of a previously unknown Y-chromosomal haplogroup was announced,
which resulted in a slight adjustment of the estimated age of the human Y-MRCA.
By definition, it is not necessary that the Y-MRCA and the mt-MRCA should have lived at the same time.
While estimates as of 2014 suggested the possibility that the two individuals may well have been roughly contemporaneous, the discovery of the archaic Y-haplogroup has pushed back the estimated age of the Y-MRCA beyond the most likely age of the mt-MRCA. As of 2015, estimates of the age of the Y-MRCA range around 200,000 to 300,000 years ago, roughly consistent with the emergence of anatomically modern humans.
Y-chromosomal data taken from a Neanderthal from El Sidrón, Spain, produced a Y-T-MRCA (time to Y-MRCA) of 588,000 years ago for Neanderthal and Homo sapiens patrilineages, dubbed ante Adam, and 275,000 years ago for Y-MRCA.
Definition
The Y-chromosomal most recent common ancestor is the most recent common ancestor of the Y-chromosomes found in currently living human males.
Due to the definition via the "currently living" population, the identity of a MRCA, and by extension of the human Y-MRCA, is time-dependent (it depends on the moment in time intended by the term "currently").
The MRCA of a population may move forward in time as archaic lineages within the population go extinct:
once a lineage has died out, it is irretrievably lost. This mechanism can thus only shift the title of Y-MRCA forward in time. Such an event could be due to the total extinction of several basal haplogroups.
The same holds for the concepts of matrilineal and patrilineal MRCAs: it follows from the definition of Y-MRCA that he had at least two sons who both have unbroken lineages that have survived to the present day. If the lineages of all but one of those sons die out, then the title of Y-MRCA shifts forward from the remaining son through his patrilineal descendants, until the first descendant is reached who had at least two sons who both have living, patrilineal descendants. The title of Y-MRCA is not permanently fixed to a single individual, and the Y-MRCA for any given population would himself have been part of a population which had its own, more remote, Y-MRCA.
Although the informal name "Y-chromosomal Adam" is a reference to the biblical Adam, this should not be misconstrued as implying that the bearer of the chromosome was the only human male alive during his time.
His other male contemporaries may also have descendants alive today, but not, by definition, through solely patrilineal descent; in other words, none of them have an unbroken male line of descendants (son's son's son's … son) connecting them to currently living people.
By the nature of the concept of most recent common ancestors, these estimates can only represent a terminus ante quem ("limit before which"), until the genome of the entire population has been examined (in this case, the genome of all living humans).
Age estimate
Estimates on the age of the Y-MRCA crucially depend on the most archaic known haplogroup extant in contemporary populations. , this is haplogroup A00 (discovered in 2013). Age estimates based on this published during 2014–2015 range between 160,000 and 300,000 years, compatible with the time of emergence and early dispersal of Homo sapiens.
Method
In addition to the tendency of the title of Y-MRCA to shift forward in time, the estimate of the Y-MRCA's DNA sequence, his position in the family tree, the time when he lived, and his place of origin, are all subject to future revisions.
The following events would change the estimate of who the individual designated as Y-MRCA was:
Further sampling of Y chromosomes could uncover previously unknown divergent lineages. If this happens, Y-chromosome lineages would converge on an individual who lived further back in time.
The discovery of additional deep rooting mutations in known lineages could lead to a rearrangement of the family tree.
Revision of the Y-chromosome mutation rate (see below) can change the estimate of the time when he lived.
The time when Y-MRCA lived is determined by applying a molecular clock to human Y-chromosomes. In contrast to mitochondrial DNA (mtDNA), which has a short sequence of 16,000 base pairs, and mutates frequently, the Y chromosome is significantly longer at 60 million base pairs, and has a lower mutation rate. These features of the Y chromosome have slowed down the identification of its polymorphisms; as a consequence, they have reduced the accuracy of Y-chromosome mutation rate estimates.
Methods of estimating the age of the Y-MRCA for a population of human males whose Y-chromosomes have been sequenced are based on applying the theories of molecular evolution to the Y chromosome. Unlike the autosomes, the human Y-chromosome does not recombine often with the X chromosome during meiosis, but is usually transferred intact from father to son; however, it can recombine with the X chromosome in the pseudoautosomal regions at the ends of the Y chromosome. Mutations occur periodically within the Y chromosome, and these mutations are passed on to males in subsequent generations.
These mutations can be used as markers to identify shared patrilineal relationships. Y chromosomes that share a specific mutation are referred to as haplogroups. Y chromosomes within a specific haplogroup are assumed to share a common patrilineal ancestor who was the first to carry the defining mutation. (This assumption could be mistaken, as it is possible for the same mutation to occur more than once.) A family tree of Y chromosomes can be constructed, with the mutations serving as branching points along lineages. The Y-MRCA is positioned at the root of the family tree, as the Y chromosomes of all living males are descended from his Y chromosome.
Researchers can reconstruct ancestral Y chromosome DNA sequences by reversing mutated DNA segments to their original condition. The most likely original or ancestral state of a DNA sequence is determined by comparing human DNA sequences with those of a closely related species, usually non-human primates such as chimpanzees and gorillas. By reversing known mutations in a Y-chromosome lineage, a hypothetical ancestral sequence for the MRCA, Y-chromosomal Adam, can be inferred.
Determining the Y-MRCA's DNA sequence, and the time when he lived, involves identifying the human Y-chromosome lineages that are most divergent from each other—the lineages that share the fewest mutations with each other when compared to a non-human primate sequence in a phylogenetic tree. The common ancestor of the most divergent lineages is therefore the common ancestor of all lineages.
History of estimates
Early estimates of the age for the Y-MRCA published during the 1990s ranged between roughly 200 and 300 thousand years ago (kya).
Such estimates were later substantially revised downward, as in Thomson et al. 2000, which proposed an age of about 59,000.
This date suggested that the Y-MRCA lived about 84,000 years after his female counterpart mt-MRCA (the matrilineal most recent common ancestor), who lived 150,000–200,000 years ago.
This date also meant that Y-chromosomal Adam lived at a time very close to, and possibly after, the migration from Africa which is believed to have taken place 50,000–80,000 years ago.
One explanation given for this discrepancy in the time depths of patrilineal vs. matrilineal lineages was that females have a better chance of reproducing than males due to the practice of polygyny. When a male individual has several wives, he has effectively prevented other males in the community from reproducing and passing on their Y chromosomes to subsequent generations. On the other hand, polygyny does not prevent most females in a community from passing on their mitochondrial DNA to subsequent generations. This differential reproductive success of males and females can lead to fewer male lineages relative to female lineages persisting into the future. These fewer male lineages are more sensitive to drift and would most likely coalesce on a more recent common ancestor. This would potentially explain the more recent dates associated with the Y-MRCA.
The "hyper-recent" estimate of significantly below 100 kya was again corrected upward in studies of the early 2010s, which ranged at about 120 kya to 160 kya.
This revision was due to the rearrangement of the backbone of the Y-chromosome phylogeny following the resequencing of Haplogroup A lineages.
In 2013, Francalacci et al. reported the sequencing of male-specific single-nucleotide Y-chromosome polymorphisms (MSY-SNPs) from 1204 Sardinian males, which indicated an estimate of 180,000 to 200,000 years for the common origin of all humans through paternal lineage.
Also in 2013, Poznik et al. reported the Y-MRCA to have lived between 120,000 and 156,000 years ago, based on genome sequencing of 69 men from 9 different populations.
In addition, the same study estimated the age of Mitochondrial Eve to about 99,000 and 148,000 years. As these ranges overlap for a time-range of 28,000 years (148 to 120 kya), the results of this study have been cast in terms of the possibility that "Genetic Adam and Eve may have walked on Earth at the same time" in the popular press.
The announcement by Mendez et al. of the discovery of a previously unknown lineage, haplogroup A00, in 2013, resulted in another shift in the estimate for the age of Y-chromosomal Adam. The authors estimated the split from the other haplogroups at 338,000 years ago (95% confidence interval 237–581 kya), but later Elhaik et al. (2014) dated it to between 163,900 and 260,200 years ago (95% CI), and Karmin et al. (2015) dated it to between 192,000 and 307,000 years ago (95% CI). The same study reports that non-African populations converge to a cluster of Y-MRCAs in a window close to 50 kya (out-of-Africa migration), and an additional bottleneck for non-African populations at about 10 kya, interpreted as reflecting cultural changes increasing the variance in male reproductive success (i.e. increased social stratification) in the Neolithic.
Family tree
Initial sequencing (Karafet et al., 2008) of the human Y chromosome suggested that two most basal Y-chromosome lineages were Haplogroup A and Haplogroup BT. Haplogroup A is found at low frequencies in parts of Africa, but is common among certain hunter-gatherer groups. Haplogroup BT lineages represent the majority of African Y-chromosome lineages and virtually all non-African lineages. Y-chromosomal Adam was represented as the root of these two lineages. Haplogroup A and Haplogroup BT represented the lineages of Y-chromosomal Adam himself and of one of his sons, who had a new SNP.
Cruciani et al. 2011, determined that the deepest split in the Y-chromosome tree was found between two previously reported subclades of Haplogroup A, rather than between Haplogroup A and Haplogroup BT. Later, group A00 was found, outside of the previously known tree.
The rearrangement of the Y-chromosome family tree implies that lineages classified as Haplogroup A do not necessarily form a monophyletic clade. Haplogroup A therefore refers to a collection of lineages that do not possess the markers that define Haplogroup BT, though Haplogroup A includes the most distantly related Y chromosomes.
The M91 and P97 mutations distinguish Haplogroup A from Haplogroup BT. Within Haplogroup A chromosomes, the M91 marker consists of a stretch of 8 T nucleobase units. In Haplogroup BT and chimpanzee chromosomes, this marker consists of 9 T nucleobase units. This pattern suggested that the 9T stretch of Haplogroup BT was the ancestral version and that Haplogroup A was formed by the deletion of one nucleobase. Haplogroups A1b and A1a were considered subclades of Haplogroup A as they both possessed the M91 with 8Ts.
But according to Cruciani et al. 2011, the region surrounding the M91 marker is a mutational hotspot prone to recurrent mutations. It is therefore possible that the 8T stretch of Haplogroup A may be the ancestral state of M91 and the 9T of Haplogroup BT may be the derived state that arose by an insertion of 1T. This would explain why subclades A1b and A1a-T, the deepest branches of Haplogroup A, both possess the same version of M91 with 8Ts. Furthermore, Cruciani et al. 2011 determined that the P97 marker, which is also used to identify Haplogroup A, possessed the ancestral state in Haplogroup A but the derived state in Haplogroup BT.
Likely geographic origin
As current estimates on TMRCA converge with estimates for the age of anatomically modern humans
and well predate the Out of Africa migration, geographical origin hypotheses continue to be limited to the African continent.
According to Cruciani et al. 2011, the most basal lineages have been detected in West, Northwest and Central Africa, suggesting plausibility for the Y-MRCA living in the general region of "Central-Northwest Africa".
Scozzari et al. (2012) agreed with a plausible placement in "the north-western quadrant of the African continent" for the emergence of the A1b haplogroup.
The 2013 report of haplogroup A00 found among the Mbo people of western present-day Cameroon is also compatible with this picture.
The revision of Y-chromosomal phylogeny since 2011 has affected estimates for the likely geographical origin of Y-MRCA as well as estimates on time depth. By the same reasoning, future discovery of presently-unknown archaic haplogroups in living people would again lead to such revisions. In particular, the possible presence of between 1% and 4% Neanderthal-derived DNA in Eurasian genomes implies that the (unlikely) event of a discovery of a single living Eurasian male exhibiting a Neanderthal patrilineal line would immediately push back T-MRCA ("time to MRCA") to at least twice its current estimate. However, the discovery of a Neanderthal Y-chromosome by Mendez et al. suggests the extinction of Neanderthal patrilineages, as the lineage inferred from the Neanderthal sequence is outside of the range of contemporary human genetic variation. Questions of geographical origin would become part of the debate on Neanderthal evolution from Homo erectus.
| Biology and health sciences | Homo | Biology |
247577 | https://en.wikipedia.org/wiki/Graph%20isomorphism | Graph isomorphism | In graph theory, an isomorphism of graphs G and H is a bijection between the vertex sets of G and H
such that any two vertices u and v of G are adjacent in G if and only if and are adjacent in H. This kind of bijection is commonly described as "edge-preserving bijection", in accordance with the general notion of isomorphism being a structure-preserving bijection.
If an isomorphism exists between two graphs, then the graphs are called isomorphic and denoted as . In the case when the isomorphism is a mapping of a graph onto itself, i.e., when G and H are one and the same graph, the isomorphism is called an automorphism of G.
Graph isomorphism is an equivalence relation on graphs and as such it partitions the class of all graphs into equivalence classes. A set of graphs isomorphic to each other is called an isomorphism class of graphs. The question of whether graph isomorphism can be determined in polynomial time is a major unsolved problem in computer science, known as the graph isomorphism problem.
The two graphs shown below are isomorphic, despite their different looking drawings.
Variations
In the above definition, graphs are understood to be undirected non-labeled non-weighted graphs. However, the notion of isomorphism may be applied to all other variants of the notion of graph, by adding the requirements to preserve the corresponding additional elements of structure: arc directions, edge weights, etc., with the following exception.
Isomorphism of labeled graphs
For labeled graphs, two definitions of isomorphism are in use.
Under one definition, an isomorphism is a vertex bijection which is both edge-preserving and label-preserving.
Under another definition, an isomorphism is an edge-preserving vertex bijection which preserves equivalence classes of labels, i.e., vertices with equivalent (e.g., the same) labels are mapped onto the vertices with equivalent labels and vice versa; same with edge labels.
For example, the graph with the two vertices labelled with 1 and 2 has a single automorphism under the first definition, but under the second definition there are two auto-morphisms.
The second definition is assumed in certain situations when graphs are endowed with unique labels commonly taken from the integer range 1,...,n, where n is the number of the vertices of the graph, used only to uniquely identify the vertices. In such cases two labeled graphs are sometimes said to be isomorphic if the corresponding underlying unlabeled graphs are isomorphic (otherwise the definition of isomorphism would be trivial).
Motivation
The formal notion of "isomorphism", e.g., of "graph isomorphism", captures the informal notion that some objects have "the same structure" if one ignores individual distinctions of "atomic" components of objects in question. Whenever individuality of "atomic" components (vertices and edges, for graphs) is important for correct representation of whatever is modeled by graphs, the model is refined by imposing additional restrictions on the structure, and other mathematical objects are used: digraphs, labeled graphs, colored graphs, rooted trees and so on. The isomorphism relation may also be defined for all these generalizations of graphs: the isomorphism bijection must preserve the elements of structure which define the object type in question: arcs, labels, vertex/edge colors, the root of the rooted tree, etc.
The notion of "graph isomorphism" allows us to distinguish graph properties inherent to the structures of graphs themselves from properties associated with graph representations: graph drawings, data structures for graphs, graph labelings, etc. For example, if a graph has exactly one cycle, then all graphs in its isomorphism class also have exactly one cycle. On the other hand, in the common case when the vertices of a graph are (represented by) the integers 1, 2,... N, then the expression
may be different for two isomorphic graphs.
Whitney theorem
The Whitney graph isomorphism theorem, shown by Hassler Whitney, states that two connected graphs are isomorphic if and only if their line graphs are isomorphic, with a single exception: K3, the complete graph on three vertices, and the complete bipartite graph K1,3, which are not isomorphic but both have K3 as their line graph. The Whitney graph theorem can be extended to hypergraphs.
Recognition of graph isomorphism
While graph isomorphism may be studied in a classical mathematical way, as exemplified by the Whitney theorem, it is recognized that it is a problem to be tackled with an algorithmic approach. The computational problem of determining whether two finite graphs are isomorphic is called the graph isomorphism problem.
Its practical applications include primarily cheminformatics, mathematical chemistry (identification of chemical compounds), and electronic design automation (verification of equivalence of various representations of the design of an electronic circuit).
The graph isomorphism problem is one of few standard problems in computational complexity theory belonging to NP, but not known to belong to either of its well-known (and, if P ≠ NP, disjoint) subsets: P and NP-complete. It is one of only two, out of 12 total, problems listed in whose complexity remains unresolved, the other being integer factorization. It is however known that if the problem is NP-complete then the polynomial hierarchy collapses to a finite level.
In November 2015, László Babai, a mathematician and computer scientist at the University of Chicago, claimed to have proven that the graph isomorphism problem is solvable in quasi-polynomial time. He published preliminary versions of these results in the proceedings of the 2016 Symposium on Theory of Computing, and of the 2018 International Congress of Mathematicians. In January 2017, Babai briefly retracted the quasi-polynomiality claim and stated a sub-exponential time complexity bound instead. He restored the original claim five days later. , the full journal version of Babai's paper has not yet been published.
Its generalization, the subgraph isomorphism problem, is known to be NP-complete.
The main areas of research for the problem are design of fast algorithms and theoretical investigations of its computational complexity, both for the general problem and for special classes of graphs.
The Weisfeiler Leman graph isomorphism test can be used to heuristically test for graph isomorphism. If the test fails the two input graphs are guaranteed to be non-isomorphic. If the test succeeds the graphs may or may not be isomorphic. There are generalizations of the test algorithm that are guaranteed to detect isomorphisms, however their run time is exponential.
Another well-known algorithm for graph isomorphism is the vf2 algorithm, developed by Cordella et al. in 2001. The vf2 algorithm is a depth-first search algorithm that tries to build an isomorphism between two graphs incrementally. It uses a set of feasibility rules to prune the search space, allowing it to efficiently handle graphs with thousands of nodes. The vf2 algorithm has been widely used in various applications, such as pattern recognition, computer vision, and bioinformatics. While it has a worst-case exponential time complexity, it performs well in practice for many types of graphs.
| Mathematics | Graph theory | null |
247652 | https://en.wikipedia.org/wiki/Bulldog%20bat | Bulldog bat | The bat family Noctilionidae, commonly known as bulldog bats or fishing bats, is represented by two extant species, the greater and the lesser bulldog bats, as well as at least one fossil species, Noctilio lacrimaelunaris, from the Miocene of Argentina. The naked bulldog bat (Cheiromeles torquatus) does not belong to this family, but to the family Molossidae, the free-tailed bats. They are found near water in the Neotropics, from Mexico to Argentina and also in the Caribbean islands. In these areas they can be found roosting in groups within hollow trees, caves, manmade homes, or other openings with enough space. While the two species exhibit different social and foraging behaviors both tend to return to a main roosting spot while also visiting other alternative roosting spots.
Description
The bulldog bats have orange to brown fur, and range in head-body length from 7 to 14 cm and weight of 20-75 g, which makes them quite large. They have relatively long legs, large feet (exceptionally so in the case of the greater bulldog bat), and strong claws. Their wings are long (up to 60 cm in spread) and narrow, and their ears are large, funnel shaped and pointed. Unusual among bats, they have cheek-pouches for storing food. They also have full lips divided by a fold of skin giving a 'hare lip' look which together with the cheek pouches gives them their bulldog-like appearance.
Their maxillae and premaxillae are fused for the strong support of the large upper medial incisors. Dental formula: 2/1, 1/1, 1/2, 3/3 = 28. The molars are tuberculosectorial. Unlike in other bats, the last cervical vertebra is not fused with the first thoracic. The wing second finger has a long metacarpal and a vestigial phalanx. The ischia are fused to each other and to the sacrum. The latter is keel-like.
Ecology and behaviour
The species of lesser bulldog bats are insectivorous, and while the greater bulldog bats also eat insects, their chief food is fish (piscivorous). They use their echolocation to pinpoint the ripples they make on the surfaces of water.
The greater bulldog bat trawls the water with its long, curved talons approximately 2–3 cm below the surface. It makes sweeps of between 30 cm and 3 m before ascending and turning to make a return sweep. In a single night, the bat may catch 20-30 small fish in this way.
| Biology and health sciences | Bats | Animals |
247827 | https://en.wikipedia.org/wiki/Gar | Gar | Gars are an ancient group of ray-finned fish in the family Lepisosteidae. They comprise seven living species of fish in two genera that inhabit fresh, brackish, and occasionally marine waters of eastern North America, Central America and Cuba in the Caribbean, though extinct members of the family were more widespread. They are the only surviving members of the Ginglymodi, a clade of fish which first appeared during the Triassic period, over 240 million years ago, and are one of only two surviving groups of holosteian fish, alongside the bowfins, which have a similar distribution.
Gars have elongated bodies that are heavily armored with ganoid scales, and fronted by similarly elongated jaws filled with long, sharp teeth. Gars are sometimes referred to as "garpike", but are not closely related to pike, which are in the fish family Esocidae. All of the gars are relatively large fish, but the alligator gar (Atractosteus spatula) is the largest; the alligator gar often grows to a length over and a weight over , and specimens of up to in length have been reported. Unusually, their vascularised swim bladders can function as lungs, and most gars surface periodically to take a gulp of air. Gar flesh is edible and the hard skin and scales of gars are used by humans, but gar eggs are highly toxic.
Etymology
The name "gar" was originally used for a species of needlefish (Belone belone) found in the North Atlantic and likely took its name from the Old English word for "spear". Belone belone is now more commonly referred to as the "garfish" or "gar fish" to avoid confusion with the North American gars of the family Lepisosteidae. Confusingly, the name "garfish" is also commonly used for a number of other species of the related genera Strongylura, Tylosurus, and Xenentodon of the family Belonidae.
The generic name Lepisosteus comes from the Greek lepis meaning "scale" and osteon meaning "bone". Atractosteus is similarly derived from Greek, in this case from atraktos, meaning arrow.
Evolution
Evolutionary history
Gars are considered to be the only surviving members of the Ginglymodi, a group of bony fish that flourished in the Mesozoic. The oldest known ginglymodians appeared during the Middle Triassic, over 240 million years ago. Because they have the slowest known rate of molecular evolution among all jawed vertebrates, it has also slowed down their rate of speciation. The closest living relatives of gars are the bowfin, with the gars and bowfin together forming the clade Holostei; both lineages diverged during the Late Permian.
The closest extinct relatives of gar are the Obaichthyidae, an extinct group of gar-like fishes from the Early Cretaceous of Africa and South America, which likely diverged from the ancestors of true gars during the Late Jurassic. The oldest anatomically modern gar is Nhanulepisosteus from the Upper Jurassic (Kimmeridgian) of Mexico, around 157 million years old. Nhanulepisosteus inhabited a marine environment unlike modern gars, indicating that gars may have originally been marine fish prior to invading freshwater habitats before the Early Cretaceous. Although most succeeding gar fossils are known from freshwater environments, at least some marine gars are known to have persisted into the Late Cretaceous, with the likely marine Herreraichthys known from Mexico and the definitely marine Grandemarinus known from Morocco.
Gars diversified in western North America throughout the Early Cretaceous. Atractosteus and Lepisosteus had already diverged by the end of the Early Cretaceous, about 105 million years ago. From western North America, gars dispersed to regions as disparate as Africa, India, South America and Europe, and fossil remains of gars were widespread worldwide by the end of the Cretaceous.
Several different gar genera survived the Cretaceous-Paleogene extinction event, although they remained restricted to North America and Europe after this point. One species (Atractosteus grandei, a relative of the modern alligator gar) is the oldest known articulated vertebrate specimen of the Cenozoic, with one fossil specimen dated to just a few thousand years after the Chicxulub impact, indicating a rapid recovery of freshwater ecosystems. Two short-snouted gar genera, Masillosteus and Cuneatus, are known from the Eocene in western North America and Europe, but disappear shortly afterwards. Lepisosteus and Atractosteus show a similar initial distribution and eventual contraction, but both genera dispersed to eastern North America prior to their disappearance from western North America and Europe, with Atractosteus also dispersing further south to the Neotropics. Eastern North America has since served as a vital refugium for gars, with Lepisosteus undergoing a diversification throughout it.
Phylogeny
The following phylogeny of extant and fossil gar genera was found by Brownstein et al. (2022):
A slightly different phylogeny was found by Cooper et al (2023):
Distribution
Fossils indicate that gars formerly had a wider distribution, having been found on every continent except Australia and Antarctica. Living gars are confined to North America. The distribution of the gars in North America lies mainly in the shallow, brackish waters off of Texas, Louisiana, and the eastern coast of Mexico, as well as in some of the rivers and lakes that flow into them. A few populations are also present in the Great Lakes region of the United States, living in similar shallow waters.
Anatomy
Scales
Gar bodies are elongated, heavily armored with ganoid scales, and fronted by similarly elongated jaws filled with long, sharp teeth. Their tails are heterocercal, and the dorsal fins are close to the tail.
Swim bladder
As their vascularised swim bladders can function as lungs, most gars surface periodically to take a gulp of air, doing so more frequently in stagnant or warm water when the concentration of oxygen in the water is low. Experiments on the swim bladder has shown that the temperature of the water affects which respiration method the gar will use - aerial or aquatic. They increase the aerial breathing rate (breathing air) as the temperature of the water is increased. Gars can live completely submerged in oxygenated water without access to air and remain healthy while also being able to survive in deoxygenated water if allowed access to air. This adaptation can be the result of environmental pressures and behavioral factors. As a result of this organ, they are extremely resilient and able to tolerate conditions that most other fish could not survive.
Pectoral girdle
The gar has paired pectoral fins and pelvic fins, as well as an anal fin, a caudal fin, and a dorsal fin. The bone structures within the fins are important to study as they can show homology throughout the fossil record. Specifically, the pelvic girdle resembles that of other actinopterygians while still having some of its own characteristics. Gars have a postcleithrun - which is a bone that is lateral to the scapula, but do not have postpectorals. Proximally to the postcleithrum, the supracleithrum is important as it plays a critical role in opening the gar's jaws. This structure has a unique internal coracoid lamina only present in the gar species. Near the supracleithrum is the posttemporal bone, which is significantly smaller than other actinopterygians. Gars also have no clavicle bone, although elongated plates have been observed within the area.
Morphology
All the gars are relatively large fish, but the alligator gar (Atractosteus spatula) is the largest. The largest alligator gar ever caught and officially recorded was long, weighed , and was around the girth. Even the smaller species, such as Lepisosteus oculatus, are large, commonly reaching lengths of over , and sometimes much longer.
Ecology
Gars tend to be slow-moving fish except when striking at their prey. They prefer the shallow and weedy areas of rivers, lakes, and bayous, often congregating in small groups. They are voracious predators, catching their prey in their needle-like teeth with a sideways strike of the head. They feed extensively on smaller fish and invertebrates such as crabs. Gars are found across much of the eastern portion of North America. Although gars are found primarily in freshwater habitats, several species enter brackish waters and a few, most notably Atractosteus tristoechus, are sometimes found in the sea. Some gars travel from lakes and rivers through sewers to get to ponds.
Species and identification
The gar family contains seven extant species, in two genera. This list also includes definitively known fossil taxa, common names for which are based on Grande (2010):
Family Lepisosteidae
Genus †Nhanulepisosteus Brito, Alvarado-Ortega & Meunier, 2017
Genus †Britosteus Martinelli et al 2024
Genus †Masillosteus Micklich & Kappert, 2001
Genus †Cuneatus Grande, 2010 (cuneatus gar)
Tribe Lepisosteini
Genus †Herreraichthys Alvarado-Ortega et al 2016
Genus †Grandemarinus Cooper et al 2023
Genus †Oniichthys Cavin & Brito, 2001
Genus Atractosteus Rafinesque, 1820
†Atractosteus atrox (Leidy, 1873) (Green River atrox gar)
†Atractosteus grandei Brownstein & Lyson, 2022
†Atractosteus messelensis Grande, 2010
†Atractosteus simplex (Leidy, 1873) (simplex gar)
Atractosteus spatula (Lacépède, 1803) (alligator gar)
Atractosteus tristoechus (Bloch & J. G. Schneider, 1801) (Cuban gar)
Atractosteus tropicus Gill, 1863 (tropical gar)
Genus Lepisosteus Linnaeus, 1758
†Lepisosteus bemisi Grande, 2010 (Green River longnose gar)
†Lepisosteus indicus (Woodward, 1890) (Indian gar)
Lepisosteus oculatus Winchell, 1864 (spotted gar)
Lepisosteus osseus (Linnaeus, 1758) (longnose gar)
Lepisosteus platostomus Rafinesque, 1820 (shortnose gar)
Lepisosteus platyrhincus DeKay, 1842 (Florida gar)
Alligator gar
The largest member of the gar family, the alligator gar (Atractosteus spatula), can measure up to 10 feet long and weigh over 300 pounds. Its body and snout are wide and stocky, and it was named "alligator gar" because locals often mistook it for an alligator. The species can be found in Texas, Oklahoma, Louisiana, the Mississippi River, Ohio, the Missouri river, and the southern drainages into Mexico. Its habitat consists of lakes and bays with slow currents. The gars grow rapidly when young and continue to grow at a slower rate after reaching adulthood. They are deep green or yellow in color. Recreational fishing of the alligator gar became popular due to its massive size and its meat is sold for food. Over five decades of overfishing have brought it close to extinction, and man-made dams have contributed to this loss by restricting the gar's access to the flood plain areas in which it spawns. Some U.S. states have enacted laws to combat overfishing, and reintroduction programs are being carried out in some states, such as Illinois, where human activity has extirpated the gar. Before being released, each gar must meet a length requirement to ensure that it has the best chance of survival in the wild. Some states, such as Texas, restrict the number of gar that may be caught in a day, the season in which they may be caught, and the equipment anglers may use to catch them. Some states also impose a minimum length requirement to prevent gar from being caught at too early an age. Scientists have found that the alligator gar can help maintain ecosystem balance by eating invasive species such as the Asian carp, and their success in a particular area can show scientists that area may also make a suitable habitat for other migratory species.
Florida gar
The Florida gar (Lepisosteus platyrhincus) can be found in the Ocklockonee river, Florida, and Georgia, and prefers muddy or sandy bottoms with bountiful vegetation. It is commonly confused with its cousin, the spotted gar. Uneven black spots cover its head, body, and fins. Green-brown scales run along the back of its body, and the scales on its underbelly are white or yellow. This coloration, which blends well with the gar's surroundings, allows it to ambush its prey. The Florida gar has no ganoid scales on its throat. Female Florida gars grow to lengths between 13 and 34 inches, bigger than their male counterparts.
Spotted gar
The spotted gar (Lepisosteus oculatus) is a smaller species of gar, measuring just under four feet long and weighing 15 pounds on average. Like Florida gars, female spotted gars are typically larger than male spotted gars. This gar has dark spots covering its head, body, and fins. Its body is compact, and it has a shorter snout. It prefers to live in clearer shallow water with a depth of 3–5 meters, and to surround itself in foliage. Its habitat ranges from the waters of Lake Michigan, the Lake Erie Basin, the Mississippi River System, and river drainages along the northern coast of the Gulf of Mexico from the Nueces River in Texas east to the lower Apalachicola River in Florida. It shares its habitat with the alligator gar, its main predator. These smaller gar live an average of 18 years.
Shortnose gar
The shortnose gar (Lepisosteus platostomus) is found in the Mississippi River Basin, Indiana, Wisconsin, Montana, Alabama, and Louisiana. It prefers to live in lakes, swamps, and calm pools. The shortnose gar takes its name from its snout, which is shorter and broader than that of other gar species. Like the longnose gar, it has one row of teeth. The upper jaw is longer than the rest of its head. The shortnose gar is deep green or brown in color, similar to the alligator gar. Depending on the clarity of water, spots can be present on the caudal, dorsal, and anal fins. The shortnose gar has a lifespan of 20 years, reaches up to 5 pounds in weight, and grows to lengths of 24–35 inches. It consumes more invertebrates than any other gar, and their stomachs have been found to contain higher Asian carp content than any other native North American fish.
Longnose gar
The Longnose gar (Lepisosteus osseus) has a longer, narrower, more cylindrical body, and can be distinguished from other species of gar by its snout, which is more than twice the length of the rest of its head. It can reach up to 6 feet and 8 inches in length and weigh up to 35–80 pounds. Like the shortnose gar, it has only a single row of teeth. Unlike its relatives, it enters brackish water from time to time. Females are larger and live longer than the male longnose gar. Females live 22 years, and males about half as long. There are spots on the head, dorsal, anal, and caudal fins. Depending on the water clarity, the longnose gar comes in two colors. In clear water, they're a dark deep green color. In muddy waters, it is more brown in color. Edges of the ganoid scales and in between are black. These types of gar are occasionally fished by locals, and blamed for eating other fish in the rivers. The longnose gar has a large range of territory in North America, into the Gulf of Mexico. Located in Florida, Quebec, all Great Lakes except Lake Superior, Missouri, Mississippi, Texas, and northern Mexico.
Roe
The flesh of gar is edible, but its eggs contain an ichthyotoxin, a type of protein toxin which is highly toxic to humans. The protein can be denatured when brought to a temperature of 120 degrees Celsius, but as the roe's temperature does not typically reach that level when it is cooked, even cooked roe causes severe symptoms. It was once thought that the production of the toxin in gar roe was an evolutionary adaptation to provide protection for the eggs, but bluegills and channel catfish fed gar eggs in experiments remained healthy, even though they are the natural predators of the gar eggs. Crayfish fed the roe were not immune to the toxin, and most died. The roe's toxicities to humans and crayfish may be coincidences, however, and not the result of explicit natural selection.
Significance to humans
Several species are traded as aquarium fish. The hard ganoid scales of gars are sometimes used to make jewelry whereas the tough skin is used to make such items as lamp shades. Historically, Native Americans used gar scales as arrowheads, native Caribbeans used the skin for breastplates, and early American pioneers covered the blades of their plows with gar skin. It is suspected that gars have an unusually strong DNA repair apparatus. If confirmed by further studies, it could be used in medical treatments against human diseases like cancer.
Not much is known about the precise function of the gar in Native American religion and culture other than the ritual "garfish dances" that have been performed by Creek and Chickasaw tribes.
| Biology and health sciences | Holosteans | Animals |
247844 | https://en.wikipedia.org/wiki/Midnight | Midnight | Midnight is the transition time from one day to the next – the moment when the date changes, on the local official clock time for any particular jurisdiction. By clock time, midnight is the opposite of noon, differing from it by 12 hours.
Solar midnight is the time opposite to solar noon, when the Sun is closest to the nadir, and the night is equidistant from sunset and sunrise. Due to the advent of time zones, which regularize time across a range of meridians, and daylight saving time, solar midnight rarely coincides with 12 midnight on the clock. Solar midnight depends on longitude and time of the year rather than on time zone. In ancient Roman timekeeping, midnight was halfway between dusk and dawn (i.e., solar midnight), varying according to the seasons.
In some Slavic languages, "midnight" has an additional geographic association with "north" (as "noon" does with "south"). Modern Polish, Belarusian, Ukrainian, and Serbian languages preserve this association with their words for "midnight" or "half-night" (północ, поўнач, північ, пoнoħ) also meaning "north".
Start and end of day
Midnight marks the beginning and ending of each day in civil time throughout the world. As the dividing point between one day and another, midnight defies easy classification as either part of the preceding day or of the following day. Though there is no global unanimity on the issue, most often midnight is considered the start of a new day and is associated with the hour 00:00.
Strictly speaking, it is incorrect to use "a.m." and "p.m." when referring to noon or midnight. The abbreviation a.m. stands for ante meridiem or before noon, and p.m. stands for post meridiem or after noon. Since noon is neither before nor after noon, and midnight is exactly twelve hours before and after noon, neither abbreviation is correct. However, many digital representations of time are configured to require an "a.m." or "p.m." designation, preventing the correct absence of such designators at midnight. In such cases, there is no international standard defining which arbitrary selection is best.
In the United States and Canada, digital clocks and computers commonly display 12 a.m. at midnight. The 30th edition of the U.S. Government Style Manual (2008), in sections 9.54 and 12.9b, recommended the use of "12 a.m." for midnight and "12 p.m." for noon. However, the previous 29th edition of the U.S. Government Printing Office Style Manual (2000), in section 12.9, recommended the opposite. There is no further record documenting this change. The US National Institute of Standards and Technology (NIST) recommends avoiding confusion altogether by using "11:59 pm" or "12:01 am" and the intended date instead of "midnight" or "12:00 am".
There are several common approaches to identifying and distinguishing the precise start and end of any given day.
Use of a 24-hour clock can remove ambiguity. The "midnight" term can be avoided altogether if the end of the day is noted as 24:00 and the beginning of the day as 00:00. While both notations refer to the same moment in time, the choice of notation allows its association with the previous night or with the following morning.
"Midnight" can be augmented with additional disambiguating information. A day and time of day may be explicitly identified together, for example "midnight Saturday night." Alternatively, midnight as the division between days may be highlighted by identifying the pair of days so divided: "midnight Saturday/Sunday" or "midnight December 14/15."
The approach recommended by the NIST ("11:59 p.m." or "12:01 a.m." instead of midnight) can be particularly helpful when any ambiguity can have serious consequences, such as with contracts and other legal instruments.
A clear convention may be legally defined or culturally promulgated. For example, the Hebrew calendar associates the start of a new day with sundown and midnight being a relative hour falling six hours after sundown. Similarly, in traditional Arabic time at sunset, which marked the start of each new day, clocks were reset to 12:00. As noted above, however, such conventions or definitions may not be uniformly observed.
The International Organization for Standardization (ISO) in specification ISO 8601 states: "00:00:00" may be used to refer to midnight corresponding to the instant at the beginning of a calendar day; and "24:00:00" to refer to midnight corresponding to the instant at the end of a calendar day.
The AP Stylebook assigns "midnight" to the day that is ending, not the day beginning.
| Physical sciences | Celestial mechanics | Astronomy |
247989 | https://en.wikipedia.org/wiki/Lactobacillus%20delbrueckii%20subsp.%20bulgaricus | Lactobacillus delbrueckii subsp. bulgaricus | Lactobacillus bulgaricus is the main bacterium used for the production of yogurt. It also plays a crucial role in the ripening of some cheeses, as well as in other processes involving naturally fermented products. It is defined as homofermentive lactic acid bacteria due to lactic acid being the single end product of its carbohydrate digestion. It is also considered a probiotic.
It is a gram-positive rod that may appear long and filamentous. It is non-motile and does not form spores. It is also non-pathogenic. It is regarded as aciduric or acidophilic, since it requires a low pH (around 5.4–4.6) to grow effectively. In addition, it is anaerobic. As it grows on raw dairy products, it creates and maintains the acidic environment that it needs to thrive via its production of lactic acid. In addition, it grows optimally at temperatures of 40–44 °C under anaerobic conditions. It has complex nutritional requirements which vary according to the environment. These include carbohydrates, unsaturated fatty acids, amino acids, and vitamins.
First identified in 1905 by the Bulgarian doctor Stamen Grigorov by isolating what later termed Lactobacillus Bulgaricus from a Bulgarian yogurt sample, the bacteria can be found naturally in the gastrointestinal tract of mammals living in Sofia region and along the Balkan Mountain (Stara Planina) mesoregion of Balkan peninsula. One strain, Lactobacillus bulgaricus GLB44, is extracted from the leaves of the Galanthus nivalis (snowdrop flower) in Bulgaria. The bacterium is also grown artificially in many countries.
Use
Lactobacillus delbrueckii subsp. bulgaricus is commonly used alongside Streptococcus thermophilus as a starter for making yogurt. The Lb. bulgaricus 2038 strain has been used for decades for yogurt fermentation. The two species work in synergy, with L. d. bulgaricus producing amino acids from milk proteins, which are then used by S. thermophilus. This relationship is considered to be symbiotic. Both species produce lactic acid, which gives yogurt its tart flavor and acts as a preservative. The resulting decrease in pH also partially coagulates the milk proteins, such as casein, resulting in yogurt's thickness. While fermenting milk, L. d. bulgaricus produces acetaldehyde, one of the main yogurt aroma components. Some strains of L. d. bulgaricus, such as L. bulgaricus GLB44, also produce bacteriocins, which have been shown to kill undesired bacteria in vitro. The viability of Lactobacillus delbrueckii subsp. bulgaricus is extremely important in that it is necessary for it to be efficient at fermentation and to effectively keep the food products it produces from spoiling. Freeze-drying is the preferred method of preserving the viability of the cells, but not all cells survive this process.
Due to its usefulness in natural fermentation processes, specifically in how it makes fermented food products out of cow's milk, it has great economic importance. Some of the biggest importers of the bacterium are Japan, the United States, and the European Union.
It has also been considered a contaminant of beer due to its homofermentative production of lactic acid, an off-flavor in many styles of beer. In other styles of beer, however, lactic acid bacteria can contribute to the overall appearance, aroma, taste, and/or mouthfeel, and generally produce an otherwise pleasing sourness.
History
Lactobacillus delbrueckii subsp. bulgaricus was first identified in 1905 by Stamen Grigorov, who named it Bacillus bulgaricus.
Ilya Metchnikoff, a professor at the Pasteur Institute in Paris, researched the relationship between the longevity of Bulgarians and their consumption of yogurt. He had the idea that aging is caused by putrefactive activity, or proteolysis, by microbes that produce toxic substances in the intestine.
Proteolytic bacteria such as clostridia, which are part of the normal intestinal flora, produce toxic substances including phenols, ammonia and indols by digestion of proteins. These compounds are responsible for what Metchnikoff called intestinal auto-intoxication, which, according to him, was the cause of the physical changes associated with old age. It was already known at that time that fermentation with lactic acid bacteria inhibits the deterioration of milk because of its low pH.
Metchnikoff's research also noted that rural populations in Southeastern Europe and the Russian steppes daily consume milk fermented with lactic acid bacteria and live relatively longer than other populations. Based on these data Metchnikoff proposed that consumption of fermented milk seeds the intestine with harmless lactic acid bacteria increasing intestinal acidity and suppressing the growth of proteolytic bacteria. His results were questioned after a 1920 study showed that the bacterium could not survive in the human intestines, but the idea nevertheless started the research into actually useful probiotics.
Lactobacillus bulgaricus is a constituent in VSL#3.
In 2012 it was declared India's national microbe.
Taxonomic history
In bacterial taxonomy, the basionym for L. d. bulgaricus was "Thermobacterium bulgaricum" . The entity became Lactobacillus bulgaricus in 1973 with the work of Rugosa and Hansen, and was reclassified as a subspecies under Lactobacillus delbrueckii in 1984.
Research
Quantification in cow's milk cheese via real-time polymerase chain reaction assay
In 2017, there was a study involving the development of a real-time polymerase chain reaction (qPCR) assay for quantifying Lactobacillus delbrueckii subsp. bulgaricus as well as Streptococcus thermophilus in cow's milk cheese. The goal of this study was to create a way to identify and quantify Lactobacillus delbrueckii subsp. bulgaricus and Streptococcus thermophilus, two lactic acid producing species crucial to the fermentation and ripening of cheese, in a timely manner through the use of qPCR. Two essays using lacZ gene targeting PCR primers resulted from this study and were deemed compatible with the two lactic acid bacteria (LAB) species. This allowed for the direct quantification of Lactobacillus delbrueckii subsp. bulgaricus and Streptococcus thermophilus in cheese produced from unpasteurized cow's milk.
Effects on antigenicity of milk proteins
A study in 2012 posed the question of whether or not Lactobacillus delbrueckii subsp. bulgaricus had any effect on the antigenicity of four kinds of milk proteins, being α-lactalbumin (α-LA), β-lactoglobulin (β-LG), α-casein (α-CN), and β-casein (β-CN). These proteins are the main proteins found in cow's milk and are known to have antigenic properties in humans, especially young children and infants. 2–5% of young children and infants experience cow's milk protein allergy (CMPA), which has harmful effects on their development and may even result in death. This allergy is facilitated through the antigenicity of the milk proteins, which is the ability of the proteins to trigger an immune response in the body that can result in a number of possible allergic reactions. The study was performed by simulating digestion of unfermented milk and milk that was fermented through exposure to Lactobacillus delbrueckii subsp. bulgaricus to compare their antigenicities in order to see if fermentation had any effect on the antigenicity of the proteins. The antigenicities were measured through an enzyme-linked immunosorbent assay (ELISA). The results claimed that the fermentation of cow's milk by Lactobacillus delbrueckii subsp. bulgaricus reduced the antigenicity of α-LA and β-CN. However, it also increased the antigenicity of α-CN while β-LG was not impacted.
Subcellular membrane fluidity under cold and osmotic stress
The efficiency of lactic acid bacteria cryopreservation is not consistent and may lead to cell death. Lactobacillus delbrueckii subsp. bulgaricus has adapted to defend against cold stress. The way most cells react to the cold is by changing the fluidity of the cellular membrane, but this particular bacterium has acquired different tactics to fight against cold stress. The first way to cope with the cold is to increase viscosity by taking in compounds such as disaccharides, polysaccharides, amino acids and antioxidants. The second strategy used is performed by inducing active responses during the fermentation or post-fermentation processes. By modifying these it will change the temperature, pH and medium composition. This results in specific metabolic pathways becoming active, with the synthesis of cold shock proteins.
Survival during freeze-drying processes
In 2017, a study was done to see the effects of six different substances on the growth and freeze- drying of Lactobacillus. Using Lactobacillus as starter cultures for the dairy industry depends on the number of viable and active cells. Currently, the preferred method to preserve the bacterial cells is through freeze-drying, however this also results in some strains being killed. This is due to various complications of freeze-drying, including the formation of ice crystals, loss of membrane fluidity, and the denaturation of important macromolecules. Regardless, freeze-drying has been used for decades in microbiological research as a way to store and stabilize cultures. Six substances, being sodium chloride, sorbitol, mannitol, mannose, monosodium glutamate, and betaine were tested to determine if they had any effect on the survivability of the cells after freeze-drying. Three of the six substances added had a positive effect on the growth and freeze-drying of Lactobacillus, being sodium chloride, sorbitol, and sodium glutamate. The results suggest that these substances have protective effects on Lactobacillus delbrueckii subsp. bulgaricus in small concentrations, but have little effect or even some harmful effects in higher concentrations. The optimal concentrations for sorbitol, sodium chloride and sodium glutamate for the desired protective effects were 0.15%, 0.6%, and 0.09% respectively. This was shown to increase cell viability drastically.
Immunotherapy for cancer
According to Helen Nauts from Cancer Research Institute, on a monograph reviewing the effects of bacterial infections on multiple types of cancer, Ivan Bogdanov, a Bulgarian physician, allegedly produced a vaccine consisting of lactobacillus bulgaricus and used it to treat two patients with myeloma, inducing remission in the two cases, one dying 18 months later due to influenza, and another living 45 months (survival median at the time was about 12–18 months). However, references are internal documents and conversations among hospitals; there's no mention in English medical literature. An article from a commercial site and an alleged documentary are available (in Bulgarian).
| Biology and health sciences | Gram-positive bacteria | Plants |
248071 | https://en.wikipedia.org/wiki/Solidago | Solidago | Solidago, commonly called goldenrods, is a genus of about 100 to 120 species of flowering plants in the family Asteraceae. Most are herbaceous perennial species found in open areas such as meadows, prairies, and savannas. They are mostly native to North America, including Mexico; a few species are native to South America and Eurasia. Some American species have also been introduced into Europe and other parts of the world.
Description
Solidago species are perennials growing from woody caudices or rhizomes. Their stems range from decumbent (crawling) to ascending or erect, with a range of heights going from to over a meter. Most species are unbranched, but some do display branching in the upper part of the plant. Both leaves and stems vary from glabrous (hairless) to various forms of pubescence (strigose, strigillose, hispid, stipitate-glandular or villous). In some species, the basal leaves are shed before flowering. The leaf margins are most commonly entire, but often display heavier serration. Some leaves may display trinerved venation rather than the pinnate venation usual across Asteraceae.
The flower heads are usually of the radiate type (typical daisy flower heads with distinct ray and disc florets) but sometimes discoid (with only disc florets of mixed, sterile, male and types). Only ray florets are female, others are male, hermaphroditic or entire sterile. Head involucres are campanulate to cylindric or attenuate. Floret corollas are usually yellow, but white in the ray florets of a few species (such as Solidago bicolor); they are typically hairless. Heads usually include between 2 and 35 disc florets, but in some species this may go up to 60. Filaments are inserted closer to the base of the corolla than its middle. Numerous heads are usually grouped in complex compound inflorescences where heads are arranged in multiple racemes, panicles, corymbs, or secund arrays (with florets all on the same side).
Solidago cypselae are narrowly obconic to cylindrical in shape, and they are sometimes somewhat compressed. They have eight to 10 ribs usually and are hairless or moderately hispid. The pappus is very big with barbellate bristles.
The many goldenrod species can be difficult to distinguish, due to their similar bright, golden-yellow flower heads that bloom in late summer. Propagation is by wind-disseminated seeds or by spreading underground rhizomes which can form colonies of vegetative clones of a single plant. They are mostly short-day plants and bloom in late summer and early fall. Some species produce abundant nectar when moisture is plentiful, or when the weather is warm and sunny.
The section Ptarmicoidei is sometimes treated as a separate genus Oligoneuron, and is dropped by flat-topped to rounded corymbiform flowerheads.
Taxonomy
Solidago is in the family Asteraceae (formerly known as Compositae), a diverse and widespread clade containing approximately 23,000 species and 12 tribes, which inhabit all continents except Antarctica. Within Asteraceae, Solidago is in the tribe Astereae and the subtribe Solidagininaeae.
The genus Solidago is monophyletic as indicated by morphological characters and molecular evidence. All Solidago species are herbaceous perennials, growing from approximately 2 cm to 2.5 m tall. Yellow to white, pistillate ray flowers and yellow, perfect disc florets are characteristic of Solidago inflorescences, which have a wide range of shapes. Molecular studies using nuclear rDNA have hypothesized boundaries on the genus Solidago, but there have been difficulties in parsing out evolutionary relationships at the sub-genus scale and defining which should be included and separated from Solidago.
Solidago and related taxa
Related Asteraceae genera, such as Chrysoma, Euthamia, and Oreochrysum, have been included within Solidago at one point or another, but morphological evidence has suggested otherwise. In a study comparing morphological characters of Solidago and related subgroups, the authors consider the subjectivity of classifying a genus and how to define it within broader tendencies concerning the taxonomy of North American Asteraceae. Little to no differences were observed between Solidago and the subgroups in terms of karyotype. However, external morphological characters such as habit, or the general appearance of the plant and how a suite of traits contribute to its phenotype; pappus size; and the point of freeing of stamen filaments from the corolla tube, are useful classification schemes for Solidago, since they are applied to differentiating between Asteraceae taxa. One school of Asteraceae taxonomy thought unites all taxa sharing similar floral head structure and subsequently ignores deviation from this morphology, while another places greater weight on these morphological deviations. The authors argue that the latter opinion should be applied. Since there is no theoretical foundation for relative taxonomic importance of traits, they assert that habit should be a central trait when defining taxa, and subsequently that all the subgroups considered in their study (Brachychaeta, Chrysoma, Euthamia, Oligoneuron, and Petradoria) should be segregated from Solidago.
Results from a leaf anatomy study comparing differences in mesophyll, bundle sheath extensions, and midvein structure, among others in a suite of leaf traits, are incongruent with those in an earlier study. Based on the lack of bundle sheath extensions, it is suggested that Chrysoma, Euthamia, Gundlachia, and Petradoria should be distinct taxa and outside of Solidago. However, Brachychaeta, Brintonia, Oligoneuron, Oreochrysum, and Aster ptarmicoides should be considered as components of Solidago. To summarize, the relation of Brachychaeta and Oligoneuron to Solidago is inconsistent based on these results. Both support the separation of Chrysoma, Euthamia, and Petradoria from Solidago. A study reviews the taxonomic position of Oligoneuron relative to Solidago, as based on taxonomic evidence, treats it as separate from Solidago, similarly to Kapoor & Beaudry (1966). The first molecular phylogeny based on chloroplast DNA treats Brachychaeta, Brintonia, Oligoneuron, and Oreochrysum as constituents of Solidago. Using consensus trees from ITS data, another study found support for Oligoneuron as part of Solidago, and the findings of Zhang (1996). More recently, an analysis of combined ITS and ETS data provided additional support for the inclusion of Oligoneuron as part of Solidago.
Until the 1980s, the genus Euthamia was largely considered to be a part of Solidago due to morphological similarities between species in both genera, and a history of synonymy of Solidago lanceolata and Euthamia graminifolia. As mentioned, the lack of bundle sheath extensions in Euthamia compared to Solidago, and deviations in floral morphology present evidence for separation of these taxa. A taxonomy of Euthamia as a genus was presented, providing a detailed description of distinguishing external morphological characters, such as fibrous-roots, sessile leaves, and mostly corymbiform inflorescences.
Evolutionary relationships within Solidago
Chromosome counts and advances in molecular systematics have enabled greater understanding of evolutionary relationships within Solidago. At the time a taxonomy of Solidago was published, related taxa causing contention, such as Chrysoma, Euthamia, Oligoneuron, and Petradoria, were excluded from this genus. The number of Solidago species has remained relatively stable, around 120, with approximately 80 in North America. Due to monophyletic support for the New World taxa and taxonomic difficulties with Old World taxa, the taxonomy provided in the 1990s only includes North American taxa and thus treats Solidago as non-monophyletic. Existing molecular-based phylogenies provide monophyletic support for Solidago given its inclusion of Oligoneuron.
Chromosome counts have proven to be a valuable character in Solidago taxonomy and in elucidating the cytogeographic history of the genus. Similar chromosome counts may indicate close evolutionary relationships, while different chromosome numbers may suggest distant relationships through reproductive isolation. Chromosome counts have been studied extensively in North America; all Solidago species have a base chromosome number of x=9, but the following ploidy levels have been observed: 2x, 3x, 4x, 6x, 8x, 10x, 12x, and 14x.
Though negligible differences in karyotype among Solidago and related genera were found, Solidago taxa with multiple cytotypes are more common than those with one. Although chromosome count is a useful metric for differentiating among Solidago taxa, it may be problematic due to the frequent variation in ploidy levels. Cytogeographic patterns in the Solidago gigantea complex, with tetraploids occurring in eastern North America and hexaploids in Oregon and Washington, were observed. Cytogeographic patterns are also observed in the Solidago canadensis complex: hexaploids within S. canadensis have been observed east of the Great Plains and are treated as Solidago altissima, and diploids and tetraploids occurring in the Great Plains are treated as Solidago gilvocanescens. The taxonomic status of Solidago ptarmicoides created an extensive debate due to frequency hybridization of S. ptarmicoides with members of the Ptarmicoidei section of Solidago. It was asserted that S. ptarmicoides should be united with Solidago rather than the genus Aster due to external morphological features such as similar pappus length as well as the same chromosome base (x=9). Information about chromosome number is still a crucial part of current understanding and phylogenies of Solidago.
Ecology
Goldenrod is considered a keystone species, and has been called the single most important plant for North American pollinator biodiversity. Goldenrod species are used as a food source by the larvae of many Lepidoptera species. As many as 104 species of butterflies and moths use it as a host plant for their larvae, and 42 species of bees are goldenrod specialists, visiting only goldenrod for food. Some lepidopteran larvae bore into plant tissues and form a bulbous tissue mass called a gall around it, upon which the larva then feeds. Various parasitoid wasps find these galls and lay eggs in the larvae, penetrating the bulb with their ovipositors. Woodpeckers are known to peck open the galls and eat the insects in the center.
Goldenrods have become invasive species in many parts of the world outside their native range, including China, Japan, Europe and Africa. Solidago canadensis, which was introduced as a garden plant in Central Europe, has become common in the wild, and in Germany is considered an invasive species that displaces native vegetation from its natural habitat.
Use and cultivation
Young goldenrod leaves are edible. Traditionally, Native Americans use the seeds of some species for food. Herbal teas are sometimes made with goldenrod.
Goldenrod often is inaccurately said to cause hay fever in humans. The pollen causing this allergic reaction is produced mainly by ragweed (Ambrosia sp.), blooming at the same time as the goldenrod and pollinated by wind. Goldenrod pollen is too heavy and sticky to be blown far from the flowers, and is pollinated mainly by insects. Frequent handling of goldenrod and other flowers, however, can cause allergic reactions, sometimes irritating enough to force florists to change occupation. Goldenrods are attractive sources of nectar for bees, flies, wasps, and butterflies. Honey from goldenrods often is dark and strong because of admixtures of other nectars. However, when honey flow is strong, a light (often water-clear), spicy-tasting monofloral honey is produced. While the bees are ripening the honey produced from goldenrods, it has a rank odour and taste; the finished honey is much milder.
Goldenrods are, in some places, considered a sign of good luck or good fortune. They are considered weeds by many in North America, but they are seen as invasive plants in Europe, where British gardeners adopted goldenrod as a garden subject. Goldenrod began to gain some acceptance in U.S. gardening (other than wildflower gardening) during the 1980s.
Cultivated species
Cultivated goldenrods include S. bicolor, S. caesia, S. canadensis, S. cutleri, S. riddellii, S. rigida, S. shortii, and S. virgaurea.
A number of cultivars have been selected, including several of hybrid origin. A putative hybrid with aster, known as ×Solidaster is less unruly, with pale yellow flowers, equally suitable for dried arrangements. Molecular and other evidence points to ×Solidaster (at least the cultivar 'Lemore') being a hybrid of Solidago ptarmicoides and Solidago canadensis, the former now in Solidago, but likely the "aster" in question.
The cultivars 'Goldenmosa'
and S. × luteus 'Lemore' have gained the Royal Horticultural Society's Award of Garden Merit.
Industrial use
Inventor Thomas Edison experimented with goldenrod to produce rubber, which it contains naturally. Edison created a fertilization and cultivation process to maximize the rubber content in each plant. His experiments produced a plant that yielded as much as 12% rubber, and the new variant was named Solidago edisoni, also called Solidago edisoniana. The tires on the Model T given to him by his friend Henry Ford were made from goldenrod. Like George Washington Carver, Henry Ford was deeply interested in the regenerative properties of soil and the potential of alternative crops such as peanuts and soybeans to produce plastics, paint, fuel and other products.
Ford had long believed that the world would eventually need a substitute for gasoline, and supported the production of ethanol (or grain alcohol) as an alternative fuel. In 1942, he would showcase a car with a lightweight plastic body made from soybeans. Ford and Carver began corresponding via letter in 1934, and their mutual admiration deepened after George Washington Carver made a visit to Michigan in 1937. As Douglas Brinkley writes in Wheels for the World, his history of Ford, the automaker donated generously to the Tuskegee Institute, helping finance Carver's experiments, and Carver in turn spent a period of time helping to oversee crops at the Ford plantation in Ways, Georgia.
By the time World War II began, Ford had made repeated journeys to Tuskegee to convince George Washington Carver to come to Dearborn and help him develop a synthetic rubber to help compensate for wartime rubber shortages. Carver arrived on July 19, 1942, and set up a laboratory in an old water works building in Dearborn. He and Ford experimented with different crops, including sweet potatoes and dandelions, eventually devising a way to make the rubber substitute from goldenrod, a plant weed commercially viable. Carver died in January 1943, Ford in April 1947, but the relationship between their two institutions continued to flourish: As recently as the late 1990s, Ford awarded grants of $4 million over two years to the George Washington Carver School at Tuskegee.
Extensive process development was conducted during World War II to commercialize goldenrod as a source of rubber. The rubber is only contained in the leaves, not the stems or blooms. Typical rubber content of the leaves is 7%. The resulting rubber is of low molecular weight, resulting in an excessively tacky compound with poor tensile properties.
Traditional medicine
Solidago virgaurea is used in a traditional kidney tonic by practitioners of herbal medicine to counter inflammation and irritation caused by bacterial infections or kidney stones. Goldenrod is also used in some formulas for cleansing of the kidney or bladder during a healing fast, in conjunction with potassium broth and specific juices. Some Native American cultures traditionally chew the leaves to relieve sore throats, and the roots to relieve toothaches.
Medicinal exploration
In various assessments by the European Medicines Agency with respect to Solidago virgaurea, non-clinical data shows diuretic, anti-inflammatory, antioxidant, analgesic and spasmolytic, antibacterial, antifungal, anticancer and immunomodulatory activity. However, as no single ingredient is responsible for these effects, the whole herbal preparation of Solidago inflorescences must be considered as the active ingredient.
Cultural significance
The goldenrod is the state flower of the U.S. states of Kentucky (adopted 1926) and Nebraska (adopted 1895). Solidago altissima, tall goldenrod, was named the state wildflower of South Carolina in 2003. The sweet goldenrod (Solidago odora) is the state herb of Delaware. Goldenrod was the state flower of Alabama, but it was later rejected in favor of the camellia.
Diversity
Accepted species
Source
Solidago albopilosa E.L.Braun – whitehair goldenrod
Solidago altiplanities C.E.S. Taylor & R.J.Taylor – high plains goldenrod
Solidago altissima L. – Canada goldenrod, late goldenrod
Solidago amplexicaulis Torr. & A.Gray
Solidago arenicola B.R. Keener & Kral – southern racemose goldenrod
Solidago argentinensis López Laphitz, Rita María & Semple
Solidago arguta Ait. – Atlantic goldenrod, forest goldenrod, toothed goldenrod, cut-leaf goldenrod
Solidago aurea Spreng.
Solidago auriculata Shuttlw. ex Blake – eared goldenrod, clasping goldenrod
Solidago bartramiana Fernald
Solidago bicolor L. – white goldenrod, silverrod
Solidago brachyphylla Chapman – Dixie goldenrod
Solidago brendiae Semple
Solidago buckleyi Torr. & Gray – Buckley's goldenrod
Solidago caesia L. – wreath goldenrod, axillary goldenrod, bluestem goldenrod, woodland goldenrod
Solidago calcicola (Fernald) Fernald
Solidago californica Nutt. - California goldenrod
Solidago canadensis L. – Canada goldenrod, Canadian goldenrod, common goldenrod
Solidago capulinensis Cockerell & Andrews L. – Capulin goldenrod
Solidago chilensis Meyen
Solidago compacta Turcz.
Solidago confinis A.Gray
Solidago coreana (Nakai) H.S.Pak
Solidago correllii Semple L. – Guadalupe Mountains goldenrod
Solidago curtisii Torr. & A.Gray – mountain decumbent goldenrod, Curtis' goldenrod
Solidago dahurica (Kitagawa) Kitagawa ex Juzepczuk
Solidago decurrens Loureiro
Solidago delicatula Small – elmleaf goldenrod, smooth elm-leaf goldenrod
Solidago drummondii Torr. & A.Gray. – Drummond's goldenrod
Solidago durangensis G.L.Nesom
Solidago elongata Nutt. – West Coast Canada goldenrod, Cascade Canada goldenrod
Solidago erecta Nutt. – showy goldenrod, slender goldenrod
Solidago ericamerioides G.L.Nesom
Solidago faucibus Wieboldt – gorge goldenrod
Solidago fistulosa P.Mill. – pine-barren goldenrod
Solidago flexicaulis L. – zigzag goldenrod, broadleaf goldenrod
Solidago gattingeri Chapman – Gattinger's goldenrod
Solidago gigantea Ait. – giant goldenrod, tall goldenrod, early goldenrod, smooth goldenrod
Solidago glabra Desf.
Solidago glomerata Michx. – clustered goldenrod, skunk goldenrod
Solidago guiradonis A.Gray – Guirado's goldenrod
Solidago gypsophila G.L.Nesom
Solidago hintoniorum G.L.Nesom
Solidago hispida Muhl. ex Willd. – hairy goldenrod
Solidago houghtonii Torr. & A.Gray ex A.Gray – Houghton's goldenrod
Solidago humilis Mill.
Solidago inornata Lunell
Solidago juliae G.L.Nesom – Julia's goldenrod
Solidago juncea Ait. – early goldenrod
Solidago kralii Semple – Kral's goldenrod
Solidago kuhistanica Juz.
Solidago kurilensis Juz.
Solidago lancifolia Torr. & A.Gray – lance-leaf goldenrod
Solidago latissimifolia P.Mill. – Elliott's goldenrod
Solidago leavenworthii Torr. & A.Gray – Leavenworth's goldenrod
Solidago leiocarpa DC. in DC. &. A.DC. – Cutler's alpine goldenrod
Solidago lepida DC. – western Canada goldenrod
Solidago ludoviciana (Gray) Small – Louisiana goldenrod
Solidago macrophylla Pursh – largeleaf goldenrod
Solidago macvaughii G.L.Nesom
Solidago microglossa DC.
Solidago minutissima (Makino) Kitam.
Solidago missouriensis Nutt. – Missouri goldenrod, prairie goldenrod, Tolmie's goldenrod
Solidago mollis Bartl. – velvety goldenrod, soft goldenrod, woolly goldenrod
Solidago multiradiata Ait. – Rocky Mountain goldenrod, alpine goldenrod, northern goldenrod, manyray goldenrod
Solidago nana Nutt. – baby goldenrod, dwarf goldenrod, gray goldenrod
Solidago nemoralis Ait. – gray goldenrod, dyersweed goldenrod, old-field goldenrod
Solidago nitida Torr. & A.Gray – shiny goldenrod
Solidago odora Ait. – anise-scented goldenrod, sweet goldenrod, fragrant goldenrod
Solidago ohioensis Riddell – Ohio goldenrod
Solidago orientalis G.L.Nesom
Solidago ouachitensis C.E.S.Taylor & R.J.Taylor – Ouachita Mountains goldenrod
Solidago ovata Friesner
Solidago pacifica Juzepczuk
Solidago paniculata DC.
Solidago patagonica Phil.
Solidago patula Muhl. ex Willd. – roundleaf goldenrod, roughleaf goldenrod
Solidago petiolaris Ait. – downy ragged goldenrod
Solidago perornata Lunell
Solidago pilosa Mill.
Solidago pinetorum Small – Small's goldenrod
Solidago plumosa Small – plumed goldenrod, plumose goldenrod, Yadkin River goldenrod
Solidago pringlei Fernald
Solidago procera Aiton
Solidago ptarmicoides (Torr. & A.Gray) B.Boivin – white flat-top goldenrod, upland white aster
Solidago puberula Nutt. – downy goldenrod
Solidago pulchra Small – Carolina goldenrod
Solidago radula Nutt. – western rough goldenrod
Solidago riddellii Frank ex Riddell – Riddell's goldenrod
Solidago rigida L. – rigid goldenrod, stiff-leaf goldenrod
Solidago roanensis Porter – Roan Mountain goldenrod
Solidago rugosa P.Mill. – wrinkleleaf goldenrod, rough-stemmed goldenrod
Solidago rupestris Raf. – rock goldenrod
Solidago satanica Lunell
Solidago sciaphila Steele – shadowy goldenrod
Solidago sempervirens L. – seaside goldenrod, salt-marsh goldenrod
Solidago serotina Retz.
Solidago shortii Torr. & A.Gray – Short's goldenrod
Solidago simplex Kunth : Mt. Albert goldenrod, sticky goldenrod
Solidago spathulata DC. – coast goldenrod
Solidago speciosa Nutt. – showy goldenrod, noble goldenrod
Solidago spectabilis (D.C.Eat.) A.Gray – Nevada goldenrod, basin goldenrod
Solidago sphacelata Raf. – autumn goldenrod, false goldenrod
Solidago spithamaea M.A.Curtis – Blue Ridge goldenrod, skunk goldenrod
Solidago spiraeifolia Fisch. ex Herder
Solidago squarrosa Nutt. – stout goldenrod
Solidago stricta Ait. – wand goldenrod, willow-leaf goldenrod
Solidago tarda Mack. – Atlantic goldenrod
Solidago tortifolia Ell. – twistleaf goldenrod
Solidago uliginosa Nutt. – bog goldenrod, fall goldenrod
Solidago ulmifolia Muhl. ex Willd. – elmleaf goldenrod
Solidago velutina DC. – threenerve goldenrod, velvety goldenrod
Solidago verna M.A.Curtis – springflowering goldenrod
Solidago villosicarpa LeBlond – glandular wand goldenrod, hairy-seed goldenrod
Solidago virgaurea L. – European goldenrod
Solidago vossii J.S.Pringle & Laureto – Voss's goldenrod
Solidago wrightii A.Gray – Wright's goldenrod
Solidago yokusaiana Makino
Natural hybrids
Solidago × asperula Desf. (S. rugosa × S. sempervirens)
Solidago × beaudryi Boivin (S. rugosa × S. uliginosa)
Solidago × calcicola (Fernald) Fernald – limestone goldenrod
Solidago × erskinei Boivin (S. canadensis × S. sempervirens)
Solidago × niederederi Khek (S. canadensis × S. virgaurea)
Solidago × ovata Friesner (S. sphacelata × S. ulmifolia)
Solidago × ulmicaesia Friesner (S. caesia × S. ulmifolia)
Formerly included
Numerous species formerly considered members of Solidago are now regarded as better suited to other genera, including Brintonia, Duhaldea, Euthamia, Gundlachia, Inula, Jacobaea, Leptostelma, Olearia, Psiadia, Senecio, Sphagneticola, Symphyotrichum, and Trixis.
| Biology and health sciences | Asterales | null |
248135 | https://en.wikipedia.org/wiki/Azide | Azide | In chemistry, azide (, ) is a linear, polyatomic anion with the formula and structure . It is the conjugate base of hydrazoic acid . Organic azides are organic compounds with the formula , containing the azide functional group. The dominant application of azides is as a propellant in air bags.
Preparation
Sodium azide is made industrially by the reaction of nitrous oxide, with sodium amide in liquid ammonia as solvent:
Many inorganic azides can be prepared directly or indirectly from sodium azide. For example, lead azide, used in detonators, may be prepared from the metathesis reaction between lead nitrate and sodium azide. An alternative route is direct reaction of the metal with silver azide dissolved in liquid ammonia. Some azides are produced by treating the carbonate salts with hydrazoic acid.
Bonding
Azide is isoelectronic with carbon dioxide , cyanate , nitrous oxide , nitronium ion , molecular beryllium fluoride and cyanogen fluoride FCN. Per valence bond theory, azide can be described by several resonance structures; an important one being .
Reactions
Azide salts can decompose with release of nitrogen gas. The decomposition temperatures of the alkali metal azides are: (275 °C), (355 °C), (395 °C), and (390 °C). This method is used to produce ultrapure alkali metals:
Protonation of azide salts gives toxic hydrazoic acid in the presence of strong acids:
Azide as a ligand forms numerous transition metal azide complexes. Some such compounds are shock sensitive.
Many inorganic covalent azides (e.g., fluorine azide, chlorine azide, bromine azide, iodine azide, silicon tetraazide) have been described.
The azide anion behaves as a nucleophile; it undergoes nucleophilic substitution for both aliphatic and aromatic systems. It reacts with epoxides, causing a ring-opening; it undergoes Michael-like conjugate addition to 1,4-unsaturated carbonyl compounds.
Azides can be used as precursors of the metal nitrido complexes by being induced to release , generating a metal complex in unusual oxidation states (see high-valent iron).
Redox behaviour and trend to disproportionation
Azides have an ambivalent redox behavior: they are both oxidizing and reducing, as they are easily subject to disproportionation, as illustrated by the Frost diagram of nitrogen. This diagram shows the significant energetic instability of the hydrazoic acid (or the azide ion) surrounded by two much more stable species, the ammonium ion on the left and the molecular nitrogen on the right. As seen on the Frost diagram the disproportionation reaction lowers ∆G, the Gibbs free energy of the system , where F is the Faraday constant, z the number of electrons exchanged in the redox reaction, and E the standard electrode potential). By minimizing the energy in the system, the disproportionation reaction increases its thermodynamical stability.
Destruction by oxidation by nitrite
Azides decompose with nitrite compounds such as sodium nitrite. Each elementary redox reaction is also a comproportionation reaction because two different N-species () converge to a same one (respectively ) and is favored when the solution is acidified. This is a method of destroying residual azides, prior to disposal. In the process, nitrogen gas () and nitrogen oxides ( and NO) are formed:
Azide (-⅓) (the reductant, electron donor) is oxidized in (0), nitrous oxide () (+1), or nitric oxide (NO) (+2) while nitrite (+3) (the oxidant, electron acceptor) is simultaneously reduced to the same corresponding species in each elementary redox reaction considered here above. The respective stability of the reaction products of these three comproportionation redox reactions is in the following order: , as can be verified in the Frost diagram for nitrogen.
Applications
In 2005, about 251 tons of azide-containing compounds were annually produced in the world, the main product being sodium azide.
Primary explosives and propellants
Sodium azide is the propellant in automobile airbags. It decomposes on heating to give nitrogen gas, which is used to quickly expand the air bag:
Heavy metal azides, such as lead azide, , are shock-sensitive detonators which violently decompose to the corresponding metal and nitrogen, for example:
Silver azide and barium azide are used similarly.
Some organic azides are potential rocket propellants, an example being 2-dimethylaminoethylazide (DMAZ) .
Microbial inhibitor and undesirable side effects
Sodium azide is commonly used in the laboratory as a bacteriostatic agent to avoid microbial proliferation in abiotic control experiments in which it is important to avoid microbial activity. However, it has the disadvantage to be prone to trigger unexpected and undesirable side reactions that can jeopardize the experimental results. Indeed, the azide anion is a nucleophile and a redox-active species. Being prone to disproportionation, it can behave both as an oxidizing and as a reducing agent. Therefore, it is susceptible to interfere in an unpredictable way with many substances. For example, the azide anion can oxidize pyrite () with the formation of thiosulfate (), or reduce quinone into hydroquinone. It can also reduce nitrite into nitrous oxide , and into (zerovalent iron, ZVI). Azide can also enhance the emission in soil. A proposed explanation is the stimulation of the denitrification processes because of the azide’s role in the synthesis of denitrifying enzymes. Moreover, azide also affects the absorbance and fluorescence optical properties of the dissolved organic matter (DOM) from soils. Many other interferences are reported in the literature for biochemical and biological analyses and they should be systematically identified and first rigorously tested in the laboratory before to use azide as microbial inhibitor for a given application.
Purification of molten sodium
Sodium azide is used to purify metallic sodium in laboratories handling molten sodium used as a coolant for fast-neutron reactors.
As hydrazoic acid, the protonated form of the azide anion, has a very low reduction potential E°red = -3,09 volt, and is even a stronger reductant than lithium (E°red = -3.04 volt), dry solid sodium azide can be added to molten metallic sodium (E°red = -2,71 volt) under strict anoxic conditions (e.g., in a special anaerobic glovebox with very low residual to reduce impurities still present into the sodium bath. The reaction residue is only gaseous .
As E°ox = -E°red, it gives the following series of oxidation reactions when the redox couples are presented as reductants:
(E°ox = +3,09 volt)
(E°ox = +3,04 volt)
(E°ox = +2,71 volt)
Click chemistry
The azide functional group is commonly utilized in click chemistry through copper(I)-catalyzed azide-alkyne cycloaddition (CuAAC) reactions, where copper(I) catalyzes the cycloaddition of an organoazide to a terminal alkyne, forming a triazole.
Other uses
A very damaging and illegal usage of sodium azide is its diversion by poachers as a substitute of sodium cyanide to poison some animal species by blocking the electron transport chain in the cellular respiration process.
Safety
Azides are explosophores and respiratory poisons. Sodium azide () is as toxic as sodium cyanide (NaCN) (with an oral of 27 mg/kg in rats) and can be absorbed through the skin. When sodium azide enters in contact with an acid, it produces volatile hydrazoic acid (), as toxic and volatile as hydrogen cyanide (HCN). When accidentally present in the air of a laboratory at low concentration, it can cause irritations such as nasal stuffiness, or suffocation and death at elevated concentrations.
Heavy metal azides, such as lead azide () are primary high explosives detonable when heated or shaken. Heavy-metal azides are formed when solutions of sodium azide or vapors come into contact with heavy metals (Pb, Hg…) or their salts. Heavy-metal azides can accumulate under certain circumstances, for example, in metal pipelines and on the metal components of diverse equipment (rotary evaporators, freezedrying equipment, cooling traps, water baths, waste pipes), and thus lead to violent explosions.
| Physical sciences | Nitride salts | Chemistry |
248189 | https://en.wikipedia.org/wiki/Gaia%20hypothesis | Gaia hypothesis | The Gaia hypothesis (), also known as the Gaia theory, Gaia paradigm, or the Gaia principle, proposes that living organisms interact with their inorganic surroundings on Earth to form a synergistic and self-regulating complex system that helps to maintain and perpetuate the conditions for life on the planet.
The Gaia hypothesis was formulated by the chemist James Lovelock and co-developed by the microbiologist Lynn Margulis in the 1970s. Following the suggestion by his neighbour, novelist William Golding, Lovelock named the hypothesis after Gaia, the primordial deity who personified the Earth in Greek mythology. In 2006, the Geological Society of London awarded Lovelock the Wollaston Medal in part for his work on the Gaia hypothesis.
Topics related to the hypothesis include how the biosphere and the evolution of organisms affect the stability of global temperature, salinity of seawater, atmospheric oxygen levels, the maintenance of a hydrosphere of liquid water and other environmental variables that affect the habitability of Earth.
The Gaia hypothesis was initially criticized for being teleological and against the principles of natural selection, but later refinements aligned the Gaia hypothesis with ideas from fields such as Earth system science, biogeochemistry and systems ecology. Even so, the Gaia hypothesis continues to attract criticism, and today many scientists consider it to be only weakly supported by, or at odds with, the available evidence.
Overview
Gaian hypotheses suggest that organisms co-evolve with their environment: that is, they "influence their abiotic environment, and that environment in turn influences the biota by Darwinian process". Lovelock (1995) gave evidence of this in his second book, Ages of Gaia, showing the evolution from the world of the early thermo-acido-philic and methanogenic bacteria towards the oxygen-enriched atmosphere today that supports more complex life.
A reduced version of the hypothesis has been called "influential Gaia" in the 2002 paper "Directed Evolution of the Biosphere: Biogeochemical Selection or Gaia?" by Andrei G. Lapenis, which states the biota influence certain aspects of the abiotic world, e.g. temperature and atmosphere. This is not the work of an individual but a collective of Russian scientific research that was combined into this peer-reviewed publication. It states the coevolution of life and the environment through "micro-forces" and biogeochemical processes. An example is how the activity of photosynthetic bacteria during Precambrian times completely modified the Earth atmosphere to turn it aerobic, and thus supports the evolution of life (in particular eukaryotic life).
Since barriers existed throughout the twentieth century between Russia and the rest of the world, it is only relatively recently that the early Russian scientists who introduced concepts overlapping the Gaia paradigm have become better known to the Western scientific community. These scientists include Piotr Alekseevich Kropotkin (1842–1921) (although he spent much of his professional life outside Russia), Rafail Vasil’evich Rizpolozhensky (1862 – ), Vladimir Ivanovich Vernadsky (1863–1945), and Vladimir Alexandrovich Kostitzin (1886–1963).
Biologists and Earth scientists usually view the factors that stabilize the characteristics of a period as an undirected emergent property or entelechy of the system; as each individual species pursues its own self-interest, for example, their combined actions may have counterbalancing effects on environmental change. Opponents of this view sometimes reference examples of events that resulted in dramatic change rather than stable equilibrium, such as the conversion of the Earth's atmosphere from a reducing environment to an oxygen-rich one at the end of the Archaean and the beginning of the Proterozoic periods.
Less accepted versions of the hypothesis claim that changes in the biosphere are brought about through the coordination of living organisms and maintain those conditions through homeostasis. In some versions of Gaia philosophy, all lifeforms are considered part of one single living planetary being called Gaia. In this view, the atmosphere, the seas and the terrestrial crust would be results of interventions carried out by Gaia through the coevolving diversity of living organisms.
The Gaia paradigm was an influence on the deep ecology movement.
Details
The Gaia hypothesis posits that the Earth is a self-regulating complex system involving the biosphere, the atmosphere, the hydrospheres and the pedosphere, tightly coupled as an evolving system. The hypothesis contends that this system as a whole, called Gaia, seeks a physical and chemical environment optimal for contemporary life.
Gaia evolves through a cybernetic feedback system operated by the biota, leading to broad stabilization of the conditions of habitability in a full homeostasis. Many processes in the Earth's surface, essential for the conditions of life, depend on the interaction of living forms, especially microorganisms, with inorganic elements. These processes establish a global control system that regulates Earth's surface temperature, atmosphere composition and ocean salinity, powered by the global thermodynamic disequilibrium state of the Earth system.
The existence of a planetary homeostasis influenced by living forms had been observed previously in the field of biogeochemistry, and it is being investigated also in other fields like Earth system science. The originality of the Gaia hypothesis relies on the assessment that such homeostatic balance is actively pursued with the goal of keeping the optimal conditions for life, even when terrestrial or external events menace them.
Regulation of global surface temperature
Since life started on Earth, the energy provided by the Sun has increased by 25–30%; however, the surface temperature of the planet has remained within the levels of habitability, reaching quite regular low and high margins. Lovelock has also hypothesised that methanogens produced elevated levels of methane in the early atmosphere, giving a situation similar to that found in petrochemical smog, similar in some respects to the atmosphere on Titan. This, he suggests, helped to screen out ultraviolet light until the formation of the ozone layer, maintaining a degree of homeostasis. However, the Snowball Earth research has suggested that "oxygen shocks" and reduced methane levels led, during the Huronian, Sturtian and Marinoan/Varanger Ice Ages, to a world that very nearly became a solid "snowball". These epochs are evidence against the ability of the pre Phanerozoic biosphere to fully self-regulate.
Processing of the greenhouse gas CO2, explained below, plays a critical role in the maintenance of the Earth temperature within the limits of habitability.
The CLAW hypothesis, inspired by the Gaia hypothesis, proposes a feedback loop that operates between ocean ecosystems and the Earth's climate. The hypothesis specifically proposes that particular phytoplankton that produce dimethyl sulfide are responsive to variations in climate forcing, and that these responses lead to a negative feedback loop that acts to stabilise the temperature of the Earth's atmosphere.
Currently the increase in human population and the environmental impact of their activities, such as the multiplication of greenhouse gases may cause negative feedbacks in the environment to become positive feedback. Lovelock has stated that this could bring an extremely accelerated global warming, but he has since stated the effects will likely occur more slowly.
Daisyworld simulations
In response to the criticism that the Gaia hypothesis seemingly required unrealistic group selection and cooperation between organisms, James Lovelock and Andrew Watson developed a mathematical model, Daisyworld, in which ecological competition underpinned planetary temperature regulation.
Daisyworld examines the energy budget of a planet populated by two different types of plants, black daisies and white daisies, which are assumed to occupy a significant portion of the surface. The colour of the daisies influences the albedo of the planet such that black daisies absorb more light and warm the planet, while white daisies reflect more light and cool the planet. The black daisies are assumed to grow and reproduce best at a lower temperature, while the white daisies are assumed to thrive best at a higher temperature. As the temperature rises closer to the value the white daisies like, the white daisies outreproduce the black daisies, leading to a larger percentage of white surface, and more sunlight is reflected, reducing the heat input and eventually cooling the planet. Conversely, as the temperature falls, the black daisies outreproduce the white daisies, absorbing more sunlight and warming the planet. The temperature will thus converge to the value at which the reproductive rates of the plants are equal.
Lovelock and Watson showed that, over a limited range of conditions, this negative feedback due to competition can stabilize the planet's temperature at a value which supports life, if the energy output of the Sun changes, while a planet without life would show wide temperature changes. The percentage of white and black daisies will continually change to keep the temperature at the value at which the plants' reproductive rates are equal, allowing both life forms to thrive.
It has been suggested that the results were predictable because Lovelock and Watson selected examples that produced the responses they desired.
Regulation of oceanic salinity
Ocean salinity has been constant at about 3.5% for a very long time. Salinity stability in oceanic environments is important as most cells require a rather constant salinity and do not generally tolerate values above 5%. The constant ocean salinity was a long-standing mystery, because no process counterbalancing the salt influx from rivers was known. Recently it was suggested that salinity may also be strongly influenced by seawater circulation through hot basaltic rocks, and emerging as hot water vents on mid-ocean ridges. However, the composition of seawater is far from equilibrium, and it is difficult to explain this fact without the influence of organic processes. One suggested explanation lies in the formation of salt plains throughout Earth's history. It is hypothesized that these are created by bacterial colonies that fix ions and heavy metals during their life processes.
In the biogeochemical processes of Earth, sources and sinks are the movement of elements. The composition of salt ions within our oceans and seas is: sodium (Na+), chlorine (Cl−), sulfate (SO42−), magnesium (Mg2+), calcium (Ca2+) and potassium (K+). The elements that comprise salinity do not readily change and are a conservative property of seawater. There are many mechanisms that change salinity from a particulate form to a dissolved form and back. Considering the metallic composition of iron sources across a multifaceted grid of thermomagnetic design, not only would the movement of elements hypothetically help restructure the movement of ions, electrons, and the like, but would also potentially and inexplicably assist in balancing the magnetic bodies of the Earth's geomagnetic field. The known sources of sodium i.e. salts are when weathering, erosion, and dissolution of rocks are transported into rivers and deposited into the oceans.
The Mediterranean Sea as being Gaia's kidney is found (here) by Kenneth J. Hsu, a correspondence author in 2001. Hsu suggests the "desiccation" of the Mediterranean is evidence of a functioning Gaia "kidney". In this and earlier suggested cases, it is plate movements and physics, not biology, which performs the regulation. Earlier "kidney functions" were performed during the "deposition of the Cretaceous (South Atlantic), Jurassic (Gulf of Mexico), Permo-Triassic (Europe), Devonian (Canada), and Cambrian/Precambrian (Gondwana) saline giants."
Regulation of oxygen in the atmosphere
The Gaia hypothesis states that the Earth's atmospheric composition is kept at a dynamically steady state by the presence of life. The atmospheric composition provides the conditions that contemporary life has adapted to. All the atmospheric gases other than noble gases present in the atmosphere are either made by organisms or processed by them.
The stability of the atmosphere in Earth is not a consequence of chemical equilibrium. Oxygen is a reactive compound, and should eventually combine with gases and minerals of the Earth's atmosphere and crust. Oxygen only began to persist in the atmosphere in small quantities about 50 million years before the start of the Great Oxygenation Event. Since the start of the Cambrian period, atmospheric oxygen concentrations have fluctuated between 15% and 40% of atmospheric volume. Traces of methane (at an amount of 100,000 tonnes produced per year) should not exist, as methane is combustible in an oxygen atmosphere.
Dry air in the atmosphere of Earth contains roughly (by volume) 78.09% nitrogen, 20.95% oxygen, 0.93% argon, 0.039% carbon dioxide, and small amounts of other gases including methane. Lovelock originally speculated that concentrations of oxygen above about 25% would increase the frequency of wildfires and conflagration of forests. This mechanism, however, would not raise oxygen levels if they became too low. If plants can be shown to robustly over-produce O2 then perhaps only the high oxygen forest fires regulator is necessary. Recent work on the findings of fire-caused charcoal in Carboniferous and Cretaceous coal measures, in geologic periods when O2 did exceed 25%, has supported Lovelock's contention.
Processing of CO2
Gaia scientists see the participation of living organisms in the carbon cycle as one of the complex processes that maintain conditions suitable for life. The only significant natural source of atmospheric carbon dioxide (CO2) is volcanic activity, while the only significant removal is through the precipitation of carbonate rocks. Carbon precipitation, solution and fixation are influenced by the bacteria and plant roots in soils, where they improve gaseous circulation, or in coral reefs, where calcium carbonate is deposited as a solid on the sea floor. Calcium carbonate is used by living organisms to manufacture carbonaceous tests and shells. Once dead, the living organisms' shells fall. Some arrive at the bottom of shallow seas where the heat and pressure of burial, and/or the forces of plate tectonics, eventually convert them to deposits of chalk and limestone. Much of the falling dead shells, however, redissolve into the ocean below the carbon compensation depth.
One of these organisms is Emiliania huxleyi, an abundant coccolithophore algae which may have a role in the formation of clouds. CO2 excess is compensated by an increase of coccolithophorid life, increasing the amount of CO2 locked in the ocean floor. Coccolithophorids, if the CLAW Hypothesis turns out to be supported (see "Regulation of Global Surface Temperature" above), could help increase the cloud cover, hence control the surface temperature, help cool the whole planet and favor precipitation necessary for terrestrial plants. Lately the atmospheric CO2 concentration has increased and there is some evidence that concentrations of ocean algal blooms are also increasing.
Lichen and other organisms accelerate the weathering of rocks in the surface, while the decomposition of rocks also happens faster in the soil, thanks to the activity of roots, fungi, bacteria and subterranean animals. The flow of carbon dioxide from the atmosphere to the soil is therefore regulated with the help of living organisms. When CO2 levels rise in the atmosphere the temperature increases and plants grow. This growth brings higher consumption of CO2 by the plants, who process it into the soil, removing it from the atmosphere.
History
Precedents
The idea of the Earth as an integrated whole, a living being, has a long tradition. The mythical Gaia was the primal Greek goddess personifying the Earth, the Greek version of "Mother Nature" (from Ge = Earth, and Aia = PIE grandmother), or the Earth Mother. James Lovelock gave this name to his hypothesis after a suggestion from the novelist William Golding, who was living in the same village as Lovelock at the time (Bowerchalke, Wiltshire, UK). Golding's advice was based on Gea, an alternative spelling for the name of the Greek goddess, which is used as prefix in geology, geophysics and geochemistry. Golding later made reference to Gaia in his Nobel prize acceptance speech.
In the eighteenth century, as geology consolidated as a modern science, James Hutton maintained that geological and biological processes are interlinked. Later, the naturalist and explorer Alexander von Humboldt recognized the coevolution of living organisms, climate, and Earth's crust. In the twentieth century, Vladimir Vernadsky formulated a theory of Earth's development that is now one of the foundations of ecology. Vernadsky was a Ukrainian geochemist and was one of the first scientists to recognize that the oxygen, nitrogen, and carbon dioxide in the Earth's atmosphere result from biological processes. During the 1920s he published works arguing that living organisms could reshape the planet as surely as any physical force. Vernadsky was a pioneer of the scientific bases for the environmental sciences. His visionary pronouncements were not widely accepted in the West, and some decades later the Gaia hypothesis received the same type of initial resistance from the scientific community.
Also in the turn to the 20th century Aldo Leopold, pioneer in the development of modern environmental ethics and in the movement for wilderness conservation, suggested a living Earth in his biocentric or holistic ethics regarding land.
Another influence for the Gaia hypothesis and the environmental movement in general came as a side effect of the Space Race between the Soviet Union and the United States of America. During the 1960s, the first humans in space could see how the Earth looked as a whole. The photograph Earthrise taken by astronaut William Anders in 1968 during the Apollo 8 mission became, through the Overview Effect an early symbol for the global ecology movement.
Formulation of the hypothesis
Lovelock started defining the idea of a self-regulating Earth controlled by the community of living organisms in September 1965, while working at the Jet Propulsion Laboratory in California on methods of detecting life on Mars. The first paper to mention it was Planetary Atmospheres: Compositional and other Changes Associated with the Presence of Life, co-authored with C.E. Giffin. A main concept was that life could be detected in a planetary scale by the chemical composition of the atmosphere. According to the data gathered by the Pic du Midi observatory, planets like Mars or Venus had atmospheres in chemical equilibrium. This difference with the Earth atmosphere was considered to be a proof that there was no life in these planets.
Lovelock formulated the Gaia Hypothesis in journal articles in 1972 and 1974, followed by a popularizing 1979 book Gaia: A new look at life on Earth. An article in the New Scientist of February 6, 1975, and a popular book length version of the hypothesis, published in 1979 as The Quest for Gaia, began to attract scientific and critical attention.
Lovelock called it first the Earth feedback hypothesis, and it was a way to explain the fact that combinations of chemicals including oxygen and methane persist in stable concentrations in the atmosphere of the Earth. Lovelock suggested detecting such combinations in other planets' atmospheres as a relatively reliable and cheap way to detect life.
Later, other relationships such as sea creatures producing sulfur and iodine in approximately the same quantities as required by land creatures emerged and helped bolster the hypothesis.
In 1971 microbiologist Dr. Lynn Margulis joined Lovelock in the effort of fleshing out the initial hypothesis into scientifically proven concepts, contributing her knowledge about how microbes affect the atmosphere and the different layers in the surface of the planet. The American biologist had also awakened criticism from the scientific community with her advocacy of the theory on the origin of eukaryotic organelles and her contributions to the endosymbiotic theory, nowadays accepted. Margulis dedicated the last of eight chapters in her book, The Symbiotic Planet, to Gaia. However, she objected to the widespread personification of Gaia and stressed that Gaia is "not an organism", but "an emergent property of interaction among organisms". She defined Gaia as "the series of interacting ecosystems that compose a single huge ecosystem at the Earth's surface. Period". The book's most memorable "slogan" was actually quipped by a student of Margulis'.
James Lovelock called his first proposal the Gaia hypothesis but has also used the term Gaia theory. Lovelock states that the initial formulation was based on observation, but still lacked a scientific explanation. The Gaia hypothesis has since been supported by a number of scientific experiments and provided a number of useful predictions.
First Gaia conference
In 1985, the first public symposium on the Gaia hypothesis, Is The Earth a Living Organism? was held at University of Massachusetts Amherst, August 1–6. The principal sponsor was the National Audubon Society. Speakers included James Lovelock, Lynn Margulis, George Wald, Mary Catherine Bateson, Lewis Thomas, Thomas Berry, David Abram, John Todd, Donald Michael, Christopher Bird, Michael Cohen, and William Fields. Some 500 people attended.
Second Gaia conference
In 1988, climatologist Stephen Schneider organised a conference of the American Geophysical Union. The first Chapman Conference on Gaia, was held in San Diego, California, on March 7, 1988.
During the "philosophical foundations" session of the conference, David Abram spoke on the influence of metaphor in science, and of the Gaia hypothesis as offering a new and potentially game-changing metaphorics, while James Kirchner criticised the Gaia hypothesis for its imprecision. Kirchner claimed that Lovelock and Margulis had not presented one Gaia hypothesis, but four:
CoEvolutionary Gaia: that life and the environment had evolved in a coupled way. Kirchner claimed that this was already accepted scientifically and was not new.
Homeostatic Gaia: that life maintained the stability of the natural environment, and that this stability enabled life to continue to exist.
Geophysical Gaia: that the Gaia hypothesis generated interest in geophysical cycles and therefore led to interesting new research in terrestrial geophysical dynamics.
Optimising Gaia: that Gaia shaped the planet in a way that made it an optimal environment for life as a whole. Kirchner claimed that this was not testable and therefore was not scientific.
Of Homeostatic Gaia, Kirchner recognised two alternatives. "Weak Gaia" asserted that life tends to make the environment stable for the flourishing of all life. "Strong Gaia" according to Kirchner, asserted that life tends to make the environment stable, to enable the flourishing of all life. Strong Gaia, Kirchner claimed, was untestable and therefore not scientific.
Lovelock and other Gaia-supporting scientists, however, did attempt to disprove the claim that the hypothesis is not scientific because it is impossible to test it by controlled experiment. For example, against the charge that Gaia was teleological, Lovelock and Andrew Watson offered the Daisyworld Model (and its modifications, above) as evidence against most of these criticisms. Lovelock said that the Daisyworld model "demonstrates that self-regulation of the global environment can emerge from competition amongst types of life altering their local environment in different ways".
Lovelock was careful to present a version of the Gaia hypothesis that had no claim that Gaia intentionally or consciously maintained the complex balance in her environment that life needed to survive. It would appear that the claim that Gaia acts "intentionally" was a statement in his popular initial book and was not meant to be taken literally. This new statement of the Gaia hypothesis was more acceptable to the scientific community. Most accusations of teleologism ceased, following this conference.
Third Gaia conference
By the time of the 2nd Chapman Conference on the Gaia Hypothesis, held at Valencia, Spain, on 23 June 2000, the situation had changed significantly. Rather than a discussion of the Gaian teleological views, or "types" of Gaia hypotheses, the focus was upon the specific mechanisms by which basic short term homeostasis was maintained within a framework of significant evolutionary long term structural change.
The major questions were:
"How has the global biogeochemical/climate system called Gaia changed in time? What is its history? Can Gaia maintain stability of the system at one time scale but still undergo vectorial change at longer time scales? How can the geologic record be used to examine these questions?"
"What is the structure of Gaia? Are the feedbacks sufficiently strong to influence the evolution of climate? Are there parts of the system determined pragmatically by whatever disciplinary study is being undertaken at any given time or are there a set of parts that should be taken as most true for understanding Gaia as containing evolving organisms over time? What are the feedbacks among these different parts of the Gaian system, and what does the near closure of matter mean for the structure of Gaia as a global ecosystem and for the productivity of life?"
"How do models of Gaian processes and phenomena relate to reality and how do they help address and understand Gaia? How do results from Daisyworld transfer to the real world? What are the main candidates for "daisies"? Does it matter for Gaia theory whether we find daisies or not? How should we be searching for daisies, and should we intensify the search? How can Gaian mechanisms be collaborated with using process models or global models of the climate system that include the biota and allow for chemical cycling?"
In 1997, Tyler Volk argued that a Gaian system is almost inevitably produced as a result of an evolution towards far-from-equilibrium homeostatic states that maximise entropy production, and Axel Kleidon (2004) agreed stating: "...homeostatic behavior can emerge from a state of MEP associated with the planetary albedo"; "...the resulting behavior of a symbiotic Earth at a state of MEP may well lead to near-homeostatic behavior of the Earth system on long time scales, as stated by the Gaia hypothesis". M. Staley (2002) has similarly proposed "...an alternative form of Gaia theory based on more traditional Darwinian principles... In [this] new approach, environmental regulation is a consequence of population dynamics. The role of selection is to favor organisms that are best adapted to prevailing environmental conditions. However, the environment is not a static backdrop for evolution, but is heavily influenced by the presence of living organisms. The resulting co-evolving dynamical process eventually leads to the convergence of equilibrium and optimal conditions".
Fourth Gaia conference
A fourth international conference on the Gaia hypothesis, sponsored by the Northern Virginia Regional Park Authority and others, was held in October 2006 at the Arlington, Virginia campus of George Mason University.
Martin Ogle, Chief Naturalist, for NVRPA, and long-time Gaia hypothesis proponent, organized the event. Lynn Margulis, Distinguished University Professor in the Department of Geosciences, University of Massachusetts-Amherst, and long-time advocate of the Gaia hypothesis, was a keynote speaker. Among many other speakers: Tyler Volk, co-director of the Program in Earth and Environmental Science at New York University; Dr. Donald Aitken, Principal of Donald Aitken Associates; Dr. Thomas Lovejoy, President of the Heinz Center for Science, Economics and the Environment; Robert Corell, Senior Fellow, Atmospheric Policy Program, American Meteorological Society and noted environmental ethicist, J. Baird Callicott.
Criticism
After initially receiving little attention from scientists (from 1969 until 1977), thereafter for a period the initial Gaia hypothesis was criticized by a number of scientists, including Ford Doolittle, Richard Dawkins and Stephen Jay Gould. Lovelock has said that because his hypothesis is named after a Greek goddess, and championed by many non-scientists, the Gaia hypothesis was interpreted as a neo-Pagan religion. Many scientists in particular also criticized the approach taken in his popular book Gaia, a New Look at Life on Earth for being teleological—a belief that things are purposeful and aimed towards a goal. Responding to this critique in 1990, Lovelock stated, "Nowhere in our writings do we express the idea that planetary self-regulation is purposeful, or involves foresight or planning by the biota".
Stephen Jay Gould criticized Gaia as being "a metaphor, not a mechanism." He wanted to know the actual mechanisms by which self-regulating homeostasis was achieved. In his defense of Gaia, David Abram argues that Gould overlooked the fact that "mechanism", itself, is a metaphor—albeit an exceedingly common and often unrecognized metaphor—one which leads us to consider natural and living systems as though they were machines organized and built from outside (rather than as autopoietic or self-organizing phenomena). Mechanical metaphors, according to Abram, lead us to overlook the active or agentic quality of living entities, while the organismic metaphors of the Gaia hypothesis accentuate the active agency of both the biota and the biosphere as a whole. With regard to causality in Gaia, Lovelock argues that no single mechanism is responsible, that the connections between the various known mechanisms may never be known, that this is accepted in other fields of biology and ecology as a matter of course, and that specific hostility is reserved for his own hypothesis for other reasons.
Aside from clarifying his language and understanding of what is meant by a life form, Lovelock himself ascribes most of the criticism to a lack of understanding of non-linear mathematics by his critics, and a linearizing form of greedy reductionism in which all events have to be immediately ascribed to specific causes before the fact. He also states that most of his critics are biologists but that his hypothesis includes experiments in fields outside biology, and that some self-regulating phenomena may not be mathematically explainable.
Natural selection and evolution
Lovelock has suggested that global biological feedback mechanisms could evolve by natural selection, stating that organisms that improve their environment for their survival do better than those that damage their environment. However, in the early 1980s, W. Ford Doolittle and Richard Dawkins separately argued against this aspect of Gaia. Doolittle argued that nothing in the genome of individual organisms could provide the feedback mechanisms proposed by Lovelock, and therefore the Gaia hypothesis proposed no plausible mechanism and was unscientific. Dawkins meanwhile stated that for organisms to act in concert would require foresight and planning, which is contrary to the current scientific understanding of evolution. Like Doolittle, he also rejected the possibility that feedback loops could stabilize the system.
Margulis argued in 1999 that "Darwin's grand vision was not wrong, only incomplete. In accentuating the direct competition between individuals for resources as the primary selection mechanism, Darwin (and especially his followers) created the impression that the environment was simply a static arena". She wrote that the composition of the Earth's atmosphere, hydrosphere, and lithosphere are regulated around "set points" as in homeostasis, but those set points change with time.
Evolutionary biologist W. D. Hamilton called the concept of Gaia Copernican, adding that it would take another Newton to explain how Gaian self-regulation takes place through Darwinian natural selection. More recently Ford Doolittle building on his and Inkpen's ITSNTS (It's The Song Not The Singer) proposal proposed that differential persistence can play a similar role to differential reproduction in evolution by natural selections, thereby providing a possible reconciliation between the theory of natural selection and the Gaia hypothesis.
Criticism in the 21st century
The Gaia hypothesis continues to be broadly skeptically received by the scientific community. For instance, arguments both for and against it were laid out in the journal Climatic Change in 2002 and 2003. A significant argument raised against it are the many examples where life has had a detrimental or destabilising effect on the environment rather than acting to regulate it. Several recent books have criticised the Gaia hypothesis, expressing views ranging from "... the Gaia hypothesis lacks unambiguous observational support and has significant theoretical difficulties" to "Suspended uncomfortably between tainted metaphor, fact, and false science, I prefer to leave Gaia firmly in the background" to "The Gaia hypothesis is supported neither by evolutionary theory nor by the empirical evidence of the geological record". The CLAW hypothesis, initially suggested as a potential example of direct Gaian feedback, has subsequently been found to be less credible as understanding of cloud condensation nuclei has improved. In 2009 the Medea hypothesis was proposed: that life has highly detrimental (biocidal) impacts on planetary conditions, in direct opposition to the Gaia hypothesis.
In a 2013 book-length evaluation of the Gaia hypothesis considering modern evidence from across the various relevant disciplines, Toby Tyrrell concluded that: "I believe Gaia is a dead end. Its study has, however, generated many new and thought provoking questions. While rejecting Gaia, we can at the same time appreciate Lovelock's originality and breadth of vision, and recognize that his audacious concept has helped to stimulate many new ideas about the Earth, and to champion a holistic approach to studying it". Elsewhere he presents his conclusion "The Gaia hypothesis is not an accurate picture of how our world works". This statement needs to be understood as referring to the "strong" and "moderate" forms of Gaia—that the biota obeys a principle that works to make Earth optimal (strength 5) or favourable for life (strength 4) or that it works as a homeostatic mechanism (strength 3). The latter is the "weakest" form of Gaia that Lovelock has advocated. Tyrrell rejects it. However, he finds that the two weaker forms of Gaia—Coeveolutionary Gaia and Influential Gaia, which assert that there are close links between the evolution of life and the environment and that biology affects the physical and chemical environment—are both credible, but that it is not useful to use the term "Gaia" in this sense and that those two forms were already accepted and explained by the processes of natural selection and adaptation.
Anthropic principle
As emphasized by multiple critics, no plausible mechanism exists that would drive the evolution of negative feedback loops leading to planetary self-regulation of the climate. Indeed, multiple incidents in Earth's history (see the Medea hypothesis) have shown that the Earth and the biosphere can enter self-destructive positive feedback loops that lead to mass extinction events.
For example, the Snowball Earth glaciations appeared to result from the development of photosynthesis during a period when the Sun was cooler than it is now. These mechanisms will have some effect, but any understanding of glacial-interglacial cycles requires study of the variations in the Earth’s orbit around the Sun, the tilt of its axis of rotation, and the ‘wobble’ in that rotational movement which causes the periodicity in Northern Hemisphere insolation, thereby setting the Earth’s thermal regime. Including studies from the fields of mathematics and Earth science, the fields of geology and geography provide insight into the causes of ice ages.
Meanwhile, the removal of carbon dioxide from the atmosphere, along with the oxidation of atmospheric methane by the released oxygen, resulted in a dramatic diminishment of the greenhouse effect. The resulting expansion of the polar ice sheets decreased the overall fraction of sunlight absorbed by the Earth, resulting in a runaway ice–albedo positive feedback loop ultimately resulting in glaciation over nearly the entire surface of the Earth. However, volcanic processes at this scale should be understood as relating to the pressure exerted on the Earth’s crust, and released during periods of ice sheet retreat. Breaking out of the Earth from the frozen condition appears to have directly been due to the release of carbon dioxide and methane by volcanos, although release of methane by microbes trapped underneath the ice could also have played a part. Lesser contributions to warming would come from the fact that coverage of the Earth by ice sheets largely inhibited photosynthesis and lessened the removal of carbon dioxide from the atmosphere by the weathering of siliceous rocks. However, in the absence of tectonic activity, the snowball condition could have persisted indefinitely.
Geologic events with amplifying positive feedbacks (along with some possible biologic participation) led to the greatest mass extinction event on record, the Permian–Triassic extinction event about 250 million years ago. The precipitating event appears to have been volcanic eruptions in the Siberian Traps, a hilly region of flood basalts in Siberia. These eruptions released high levels of carbon dioxide and sulfur dioxide which elevated world temperatures and acidified the oceans. Estimates of the rise in carbon dioxide levels range widely, from as little as a two-fold increase, to as much as a twenty-fold increase. Amplifying feedbacks increased the warming to considerably greater than that to be expected merely from the greenhouse effect of carbon dioxide: these include the ice albedo feedback, the increased evaporation of water vapor (another greenhouse gas) into the atmosphere, the release of methane from the warming of methane hydrate deposits buried under the permafrost and beneath continental shelf sediments, and increased wildfires. The rising carbon dioxide acidified the oceans, leading to widespread die-off of creatures with calcium carbonate shells, killing mollusks and crustaceans like crabs and lobsters and destroying coral reefs. Their demise led to disruption of the entire oceanic food chain. It has been argued that rising temperatures may have led to disruption of the chemocline separating sulfidic deep waters from oxygenated surface waters, which led to massive release of toxic hydrogen sulfide (produced by anerobic bacteria) to the surface ocean and even into atmosphere, contributing to the (primarily methane-driven) collapse of the ozone layer, and helping to explain the die-off of terrestrial animal and plant life.
According to the weak anthropic principle, our observation of such stabilizing feedback loops is an observer selection effect. In all the universe, it is only planets with Gaian properties that could have evolved intelligent, self-aware organisms capable of asking such questions. One can imagine innumerable worlds where life evolved with different biochemistries or where the worlds had different geophysical properties such that the worlds are presently dead due to runaway greenhouse effect, or else are in perpetual Snowball, or else due to one factor or another, life has been inhibited from evolving beyond the microbial level.
If no means exists for natural selection to operate at the biosphere level, then it would appear that the anthropic principle provides the only explanation for the survival of Earth's biosphere over geologic time. But in recent years, this strictly reductionistic view has been modified by recognition that natural selection can operate at multiple levels of the biological hierarchy — not just at the level of individual organisms. Traditional Darwinian natural selection requires reproducing entities that display inheritable properties or abilities that result in their having more offspring than their competitors. Successful biospheres clearly cannot reproduce to spawn copies of themselves, and so traditional Darwinian natural selection cannot operate. A mechanism for biosphere-level selection was proposed by Ford Doolittle: Although he had been a strong and early critic of the Gaia hypothesis, he had by 2015 started to think of ways whereby Gaia might be "Darwinised", seeking means whereby the planet could have evolved biosphere-level adaptations. Doolittle has suggested that differential persistence — mere survival — could be considered a legitimate mechanism for natural selection. As the Earth passes through various challenges, the phenomenon of differential persistence enables selected entities to achieve fixation by surviving the death of their competitors. Although Earth's biosphere is not competing against other biospheres on other planets, there are many competitors for survival on this planet. Collectively, Gaia constitutes the single clade of all living survivors descended from life’s last universal common ancestor (LUCA). Various other proposals for biosphere-level selection include sequential selection, entropic hierarchy, and considering Gaia as a holobiont-like system. Ultimately speaking, differential persistence and sequential selection are variants of the anthropic principle, while entropic hierarchy and holobiont arguments may possibly allow understanding the emergence of Gaia without anthropic arguments.
| Physical sciences | Earth science basics: General | Earth science |
248234 | https://en.wikipedia.org/wiki/Booby%20trap | Booby trap | A booby trap is a device or setup that is intended to kill, harm or surprise a human or another animal. It is triggered by the presence or actions of the victim and sometimes has some form of bait designed to lure the victim towards it. The trap may be set to act upon trespassers that enter restricted areas, and it can be triggered when the victim performs an action (e.g., opening a door, picking something up, or switching something on). It can also be triggered by vehicles driving along a road, as in the case of improvised explosive devices (IEDs).
Booby traps should not be confused with mantraps which are designed to catch a person. Lethal booby traps are often used in warfare, particularly guerrilla warfare, and traps designed to cause injury or pain are also sometimes used by criminals wanting to protect drugs or other illicit property, and by some owners of legal property who wish to protect it from theft. Booby traps which merely cause discomfort or embarrassment are a popular form of practical joke.
Etymology
The Spanish word translates to "stupid, daft, naïve, simple, fool, idiot, clown, funny man, one who is easily cheated" and similar pejorative terms. The slang of , , translates to "dunce". Variations of this word exist in other languages (such as Latin), with their meaning being "to stammer".
In approximately 1590, the word began appearing in the English language as booby, meaning "stupid person, slow bird". The seabird in question was the genus Sula, with their common name being boobies. These birds have large flat feet and wide wingspans for marine habitats but are clumsy and slow on shore making them easy to catch. The birds are also known for landing aboard seagoing vessels, whereupon they have been eaten by the crew.
The phrase booby trap originally applied to schoolboy pranks, but took on its more serious connotation during World War I.
Military booby traps
A military booby trap is designed to kill or injure a person who activates its trigger, or employed to reveal the location of an enemy by setting off a signalling device. Most, but not all, military booby traps involve explosives.
Part of the skill in placing booby traps lies in exploiting natural human behaviors such as habit, self–preservation, curiosity or acquisitiveness. A common trick is to provide victims with a simple solution to a problem, for example, leaving only one door open in an otherwise secure building, luring them straight toward the firing mechanism.
An example that exploits an instinct for self–preservation was used in the Vietnam War. Spikes known as punji sticks were hidden in grassy areas. When fired upon, soldiers instinctively sought to take cover by throwing themselves down on the ground, impaling themselves on the spikes.
Many purpose–built booby–trap firing devices exist such as the highly versatile M142 universal firing device (identical to the British L5A1 or Australian F1A1), or Yugoslavian UMNOP-1 which allow a variety of different ways of triggering explosives e.g. via trip wire (either pulling it or releasing the tension on it), direct pressure on an object (e.g. standing on it), or pressure release (lift/shift something) etc.
Most explosive booby traps use between 250 g and 1 kg of explosive. Since most booby traps are rigged to detonate within a metre of the victim's body, this is adequate to kill or severely wound.
Effects
Booby traps are indiscriminate weapons. Like anti-personnel mines, they can harm civilians and noncombatants during and after the conflict. The use against civilians is prohibited by the Protocol on Mines, Booby-Traps and Other Devices, and the protocol also prohibits boobytrapping e.g. the wounded or dead, medical equipment, food, and drink.
History
A type of booby trap was referred to in an 1839 news story in The Times.
During the Vietnam War, motorcycles were rigged with explosives by the National Liberation Front and abandoned. U.S. soldiers would be tempted to ride the motorcycle and thus trigger the explosives. In addition, NLF soldiers would rig rubber band grenades and place them in huts that US soldiers would likely burn. Another popular booby trap was the "Grenade in a Can", a grenade with the safety pin removed in a container and a string attached, sometimes with the grenade's fuse mechanism modified to give a much shorter delay than the four to seven seconds typical with grenade fuses. The NLF soldiers primarily used these on doors and attached them to tripwires on jungle paths.
The CIA and Green Berets countered by booby trapping the enemy's ammunition supplies, in an operation code–named "Project Eldest Son". The propellant in a rifle or machine–gun cartridge was replaced with high explosive. Upon being fired, the sabotaged round would destroy the gun and kill or injure the shooter. Mortar shells were similarly rigged to explode when dropped down the tube, instead of launching properly. This ammunition was then carefully re–packed to eliminate any evidence of tampering, and planted in enemy munitions dumps by covert insertion teams. A sabotaged round might also be planted in a rifle magazine or machine–gun belt and left on the body of a dead NLF soldier, in anticipation that the deceased's ammo would be picked up and used by his comrades. No more than one sabotaged round would be planted in any case, magazine, or belt of ammunition, to reduce the chances of the enemy finding it no matter how diligently they inspected their supplies. False rumors and forged documents were circulated to make it appear that the Communist Chinese were supplying the NLF with defective weapons and ammunition.
Northern Ireland
During the Troubles, an ethnonationalist conflict in Northern Ireland, booby traps were used by Irish republican and Ulster loyalist paramilitaries to target British security forces and civilians. The Provisional Irish Republican Army (IRA) was the most prolific user of the booby traps during the conflict; according to the Sutton Index of Deaths, 180 people were killed during the Troubles as the result of booby trap bombs, the vast majority of them laid by the IRA. A common type of booby trap was the car bomb, which involved attaching a bomb to a car so that starting or driving it would detonate the explosive.
Middle East
Lebanese media reported on the phenomenon of intentionally concealed bombs in children's toys in 1997, citing a number of examples. Israel occupied southern Lebanon between 1982 and 2000 and during that period planted hundreds of thousands of landmines and bomblets. A report by the UK Foreign Affairs Committee in 2000 warned of the dangers of unexploded bombs in southern Lebanon, mentioning the use of "booby-trapped toys, allegedly dropped by the Israeli air force near Lebanese villages adjacent to the so-called security zone".
During the Al-Aqsa Intifada (2000—2005), some Arab–Palestinian groups made wide use of booby traps to prevent the Israeli army from entering their cities on Palestinian territories. The largest use of booby traps was in the Battle of Jenin during Operation Defensive Shield where a large number (1000–2000 bombs and booby traps according to a Palestinian militant who surrendered to Israeli forces in Jenin) of explosive devices were planted by insurgents. Booby traps had been laid in the streets of both the camp and the town, ready to be triggered if a foot snagged a tripwire or a vehicle rolled over a mine. Some of the bombs were huge, containing as much as 250 lb (110 kg) of explosives. To counter the booby traps, anti–tank and anti–personnel mines the Israeli army sent armored Caterpillar D9 bulldozers to clear the area out of any explosive device and booby trap planted. Eventually, a dozen D9 bulldozers went into action, razing the center of the refugee camp and forcing the Palestinian militants inside to surrender.
In the Israel–Hamas war, Israel's use of pagers and walkie-talkies detonations to target the Lebanese militant group Hezbollah that killed 42 people and injured 3,500 more, was condemned by some as illegal.
Gallery
Civilian use and legal ramifications
Booby traps have been applied as defensive weapons against non-military trespassers, but most jurisdictions consider the practice illegal.
Practical jokes
Instead of being used to kill, maim or injure people, booby traps can also be used for entertainment. Practical joke booby traps are typically disguised as everyday items such as cigars or packets of chewing gum, nuts or other snack items. When the victims attempts to use the item, the trap is triggered. Two of the best known examples of this are the exploding cigar and dribble glass; others include the snake nut can and shocking gum. Booby traps can also be constructed out of household or workplace items and be triggered when the victim performs a common action. Examples of this include loosening the bolts in a chair so that it collapses when sat upon, or placing a bucket of water on top of a partly open door so that when the door is fully opened, the bucket tips onto the victim.
| Technology | Military technology: General | null |
248446 | https://en.wikipedia.org/wiki/Hartsfield%E2%80%93Jackson%20Atlanta%20International%20Airport | Hartsfield–Jackson Atlanta International Airport | Hartsfield–Jackson Atlanta International Airport is the primary international airport serving Atlanta and its surrounding metropolitan area, in the U.S. state of Georgia. The airport is located south of the Downtown Atlanta district. It is named after former Atlanta mayors William B. Hartsfield and Maynard Jackson. The airport covers of land and has five parallel runways which are aligned in an east–west direction. There are three runways that are long, one runway that is long, and the longest runway at ATL measures long, which can handle the Airbus A380. Since 1998, Hartsfield–Jackson has been the world's busiest airport by passenger traffic, with the exception of 2020, when its passenger traffic dipped for that year due to travel restrictions resulting from the COVID-19 pandemic. In 2023, the airport served over 104.6 million passengers, the most of any airport in the world. Hartsfield–Jackson is also the world's busiest airport by aircraft movements. In 2024, it was again named the busiest airport in the world, and saw 2% more capacity than the year previous.
Hartsfield–Jackson is the corporate headquarters and primary hub of Delta Air Lines. With just over 1,000 flights a day to 225 domestic and international destinations, the Delta hub is the world's largest airline hub and is considered the first mega-hub in America. Additionally, Hartsfield–Jackson is the home of Delta's Technical Operations Center, which is the airline's primary maintenance, repair and overhaul arm. Aside from Delta, Hartsfield–Jackson is also an operating base for low-cost carriers Frontier Airlines, Southwest Airlines, and Spirit Airlines. The airport has international service within North America and to Latin America, Europe, Africa, Middle East and East Asia.
The airport is mostly in unincorporated areas of Clayton County, but it spills into the city limits of Atlanta, College Park, and Hapeville, in territory extending into Fulton County. The airport's domestic terminal is served by MARTA's Red and Gold rail lines.
History
Candler Field/Atlanta Municipal Airport (1925–1961)
Hartsfield–Jackson began with a five-year, rent-free lease on that was an abandoned auto racetrack named The Atlanta Speedway. The lease was signed on April 16, 1925, by Mayor Walter Sims, who committed the city to develop it into an airfield. As part of the agreement, the property was renamed Candler Field after its former owner, Coca-Cola tycoon and former Atlanta mayor Asa Candler. The first flight into Candler Field was September 15, 1926, a Florida Airways mail plane flying from Jacksonville, Florida. In May 1928, Pitcairn Aviation began service to Atlanta, followed in June 1930 by Delta Air Service. Those two airlines, later known as Eastern Air Lines and Delta Air Lines, respectively, would both use Atlanta as their chief hubs. The airport's weather station became the official location for Atlanta's weather observations on September 1, 1928, and records by the National Weather Service.
Atlanta was a busy airport from its inception, and by the end of 1930, it was third behind New York City and Chicago for regular daily flights with sixteen arriving and departing. Candler Field's first control tower opened March 1939. The March 1939 Official Aviation Guide shows fourteen weekday airline departures: ten Eastern and four Delta.
In October 1940, the U.S. government declared it a military airfield and the United States Army Air Forces operated Atlanta Army Airfield jointly with Candler Field. The Air Force used the airport primarily to service many types of transient combat aircraft. During World War II, the airport doubled in size and set a record of 1,700 takeoffs and landings in a single day, making it the nation's busiest in terms of flight operation. Atlanta Army Airfield closed after the war.
In 1942, Candler Field was renamed Atlanta Municipal Airport and by 1948, more than one million passengers passed through a war surplus hangar that served as a terminal building. Delta and Eastern had extensive networks from ATL, though Atlanta had no nonstop flights beyond Texas, St. Louis, and Chicago until 1961. Southern Airways appeared at ATL after the war and had short-haul routes around the Southeast until 1979.
In 1957, Atlanta saw its first jet airliner: a prototype Sud Aviation Caravelle that was touring the country arrived from Washington, D.C. The first scheduled turbine airliners were Capital Viscounts in June 1956; the first scheduled jets were Delta DC-8s in September 1959. The first trans-Atlantic flight was a Delta/Pan Am interchange DC-8 to Europe via Washington starting in 1964; the first scheduled international nonstops were Eastern flights to Mexico City and Jamaica in 1971–72. Nonstops to Europe started in 1978 and to Asia in 1992–93. a
Atlanta claimed to be the country's busiest airport, with more than two million passengers passing through in 1957 and, between noon and 2p.m. each day, it became the world's busiest airport. (The April 1957 OAG shows 165 weekday departures from Atlanta, including 45 between 12:05 and 2:00 PM and 20 between 2:25 and 4:25 AM.) Chicago Midway had 414-weekday departures, including 48 between 12:00 and 2:00 PM. In 1957, Atlanta was the country's ninth-busiest airline airport by flight count and about the same by passenger count.
Original Jet Terminal (1961–1980)
In late 1957, work began on a new $21 million terminal, which opened on May 3, 1961. Consisting of six pier concourses radiating from a central building, the terminal was the largest in the country and could handle over six million travelers a year; the first year, nine and a half million people passed through. In March 1962, the longest runway (9/27, now 8R) was ; runway3 was and runway 15 was long.
In 1971, the airport was named William B. Hartsfield Atlanta Airport in honor of Atlanta mayor William B. Hartsfield after his death. The name change took effect on February 28, which would have been Hartsfield's 81st birthday. The new name would be relatively brief, as it would be changed later in 1971 to William B. Hartsfield Atlanta International Airport with the growth of flights to and from Atlanta outside North America.
Midfield Terminal (1980–present)
To address the significant increase in air traffic that outstripped the capacity of the 1961 terminal, and after years of planning and design, construction began on the present midfield terminal complex in 1977 under the administration of Mayor Maynard Jackson. It was billed as the largest construction project in the South, costing $500 million. The complex was designed by Stevens & Wilkinson, Smith Hinchman & Grylls, and Minority Airport Architects & Planners. The new complex, initially consisting of the North and South Terminals, Concourses A through D, and the northern half of the present-day Concourse T (which served as the International Terminal), opened on September 21, 1980, on time and under budget. It was designed to accommodate up to 55 million passengers per year and covered 2.5 million square feet (230,000 m2). In December 1984, a fourth parallel runway was completed, and another runway was extended to the following year. To accommodate increases in international air traffic, a southern extension of Concourse T opened in 1987, and Concourse E opened in 1994 in advance of Atlanta hosting the 1996 Summer Olympics, with Concourse T subsequently being converted to use by domestic flights. MARTA rail service was extended to Hartsfield with the opening of the Airport station in 1988 (the station itself was constructed in 1979-80 as part of the terminal).
In 1999, Hartsfield–Jackson's leadership established the Development Program: "Focus On the Future," involving multiple construction projects to prepare the airport to handle a projected demand of 121 million passengers in 2015. The program was originally budgeted at $5.4 billion over ten years, but the total was revised as of 2007 to over $9 billion.
In 2001, construction of an over fifth runway (10–28) began. It was completed at the cost of $1.28 billion and opened in 2006. It bridges Interstate 285 (the Perimeter) on the airport's south side, making Hartsfield–Jackson the nation's only currently active civil airport to have a runway above an interstate (although Runway 17R/35L at Stapleton International Airport in Denver, Colorado, crossed Interstate 70 until that airport closed in 1995). The massive project, which involved putting fill dirt eleven stories high in some places, destroyed some surrounding neighborhoods and dramatically changed the scenery of Flat Rock Cemetery and Hart Cemetery, both on the airport property. It was added to help ease traffic problems caused by landing small- and mid-size aircraft on the runways used by larger planes such as the Boeing 777, which need longer runways than the smaller planes. With the fifth runway, Hartsfield–Jackson is one of only a few airports that can perform triple simultaneous landings. The fifth runway was expected to increase the capacity for landings and take-offs by 40%, from an average of 184 flights per hour to 237 flights per hour.
Along with the fifth runway, a new control tower was built to see the entire runway length. The new control tower is the tallest in the United States, over tall. The old control tower, at 231 ft, was demolished in 2006.
In 2003, the Atlanta City Council voted to rename Hartsfield–Jackson Atlanta International Airport to honor former mayor Maynard Jackson, who died four months prior. The council planned to drop Hartsfield's name from the airport, but public outcry (occurring coincidentally during a debate over the state's flag) prevented this.
In 2007, an "end-around taxiway" opened, Taxiway Victor. It is expected to save an estimated $26 million to $30 million in fuel each year by allowing airplanes landing on the northernmost runway to taxi to the gate area without preventing other aircraft from taking off. The taxiway drops about from runway elevation to allow takeoffs to continue.
After the Southeastern U.S. drought of 2007, the airport (the state's eighth-largest water user) changed to reduce water usage. This included adjusting toilets (725 commodes and 338 urinals) and 601 sinks. (The two terminals alone use a day.) It also stopped using firetrucks to spray water over aircraft when the pilot made the last landing before retirement (a water salute). The city of Macon offered to sell water to the airport through a proposed pipeline.
The Maynard H. Jackson International Terminal and Concourse F opened on the east side of the airport for international passengers in 2012. The 1980 terminal on the other end of the complex then became known as the Domestic Terminal. Prior to the opening of the International Terminal, all Atlanta-bound international passengers needed to go through TSA screening and transit to the terminal to exit the airport. The opening of the International Terminal eliminated the need for this practice, which had been in use since the opening of Concourse E in 1994.
The airport today employs about 55,300 airline, ground transportation, concessionaire, security, the federal government, the City of Atlanta, and airport tenant employees and is the largest employment center in Georgia. With a payroll of $2.4 billion, the airport has a direct and indirect economic impact of $3.2 billion on the local and regional economy and an annual regional economic impact of more than $19.8 billion.
In 2015, the airport became the first airport in the world to serve 100 million passengers in a year. The airport is routinely cited as one of the world's busiest, topping the Airports Council International rankings in 2022 and 2023.
Historical airline service
Delta and Eastern dominated the airport during the 1970s. United, Southern, Piedmont, Northwest and TWA were also present. In 1978, after airline deregulation, United no longer served Atlanta, while Southern successor Republic was the airport's third-largest carrier.
Eastern was a larger airline than Delta until deregulation in 1978, but Delta was early to adopt the hub-and-spoke route system, with Atlanta as a hub between the Midwest and Florida, giving it an advantage in the Atlanta market. When the current terminal complex opened in 1980, Delta occupied all of Concourse A and the southern side of Concourse B, while Eastern occupied the remainder of Concourse B and all of Concourse C. All other domestic airlines used Concourse D, and Concourse T (known then as the International Concourse) was used by international flights.
Eastern ceased operations in 1991. From Eastern's demise to the 1996 Summer Olympics, Delta's hub grew to occupy all of Concourse B and the southern side of Concourse T (which opened in 1987), and international flights moved to the new Concourse E (which opened in 1994). By 1996, Delta's regional affiliate Atlantic Southeast Airlines (operating as Delta Connection) relocated to the north side of Concourse C and the gates were converted for use by regional aircraft.
After Eastern ceased operation, Northwest Airlines (the successor of Republic) briefly expressed interest in establishing an Atlanta hub but ultimately decided against it. American Airlines also considered establishing an Atlanta hub around that time but decided Delta was too strong there and instead replaced Eastern's other hub in Miami. In 1992, TWA created a small hub at Atlanta and relocated to some of Eastern's former gates on Concourse C. TWA abandoned the Atlanta hub concept in 1994 leaving Delta with a monopoly hub at Atlanta.
Japan Airlines was the first Asian carrier to serve Atlanta in 1986. In December 1994, Korean Air became the second Asian carrier to serve the airport.
Atlanta-based ValuJet was established in 1993 as low-cost competition for Delta at ATL. ValuJet built up their hub on Concourse C in the following years. However, ValuJet's safety practices were questioned early, and the airline was grounded after the 1996 crash of ValuJet Flight 592. ValuJet resumed operations later that year and in 1997, it merged with AirTran Airways. AirTran would continue operating the hub and was second-largest airline at ATL through the 2000s. AirTran was acquired by Southwest Airlines in 2011, who did not serve Atlanta prior to the acquisition. AirTran was fully absorbed into Southwest in 2014, continuing to operate Atlanta as a focus city and remaining the airport's second-largest carrier. In 2024, Southwest announced it was permanently cutting 15 destinations from Atlanta, reducing its footprint from 18 gates to 11, and cutting staff.
In recent years the airport has had an increase in non-Delta flights, both due to the rapid population growth of Metro Atlanta and the airport's prominence as a major hub.
Since 2015 the airport has seen growth from low cost carriers such as Frontier Airlines and Spirit Airlines. Spirit also established Atlanta as an operating base.
In addition to the growth of the low cost carriers, international carriers have increasingly offered service to Atlanta since 2014. In 2014, Virgin Atlantic began offering direct flights to London and in 2015, the airline began offering direct flights to Manchester. In May 2016, Turkish Airlines began offering direct flights to Istanbul and Qatar Airways began Doha flights just one month later on June 1. In March 2019, WestJet began offering direct flights to Calgary, and in 2023, the airline started non-stop service to Vancouver and Winnipeg. In 2024, WestJet began non-stop service to Edmonton. Copa Airlines became the first Latin American carrier to serve the airport in December 2021 with direct flights to Panama City. In June 2022, Air Canada reintroduced Montreal service. Ethiopian Airlines started service to Atlanta in 2023, becoming the first African carrier to serve the airport since South African Airways ended service in 2006. LATAM Perú started service to Atlanta in October 2023 from Lima. Aeromexico Connect resumed service to Atlanta in January 2024 with nonstop service to Guadalajara and Monterrey. Nonstop service to Leon/Guanajuato and Mérida began in March 2024. The Mérida service ended in June. Nonstop service to Querétaro started service in August 2024. Nonstop service to Manzanillo started on November 2. Scandinavian Airlines started service to Atlanta in June 2024 with direct flights from Copenhagen. Etihad Airways will start nonstop service to Atlanta on July 2, 2025 with direct flights to Abu Dhabi.
Facilities
Terminals
Hartsfield–Jackson Atlanta International Airport has two terminals and seven concourses with a total of 192 gates. The Domestic Terminal is located on the west side of the airport and the Maynard H. Jackson Jr. International Terminal is on the east side of the airport. The Domestic Terminal has entrances on both sides, which are known as Domestic Terminal North and Domestic Terminal South. Concourse T is directly connected to the Domestic Terminal and Concourse F is directly connected to the International Terminal. The remaining five concourses (Concourses A-E) are located between the two terminals and are parallel to each other. The terminals and concourses are connected airside by the Transportation Mall, an underground pedestrian tunnel with a series of moving walkways and The Plane Train, a 24/7 underground automated people mover. A second underground walkway connecting the north sides of Concourses B and C once existed for Eastern Air Lines. Though, this underground walkway was closed in the late 2000s and is now used for the airport's baggage system.
Delta Air Lines' hub includes operations on all seven concourses. The south side of Concourse T and all of Concourses A and B are used exclusively by Delta for main line domestic flights. Delta's regional flights (operated as Delta Connection) primarily operate from the north side of Concourse C. The south side of Concourse C is used by Southwest Airlines for their operating base. All other domestic airlines operate from Concourse D or the north side of Concourse T. Some Delta and Delta Connection flights operate on Concourse D as well.
International flights operate in Concourses E and F. Concourse F is the only concourse in the airport that has a gate that can support an Airbus A380, the largest passenger aircraft in the world. All non-Delta international carriers operate their ATL flights from this terminal, including Delta's partners such as Air France, KLM, Korean Air, LATAM, Virgin Atlantic, Scandinavian, and WestJet. Aeromexico operates in Concourse E. Some WestJet flights operates in Concourse D.
International passengers arriving in Concourse F will be processed at the Customs and Border Protection checkpoint in that concourse before exiting into the landside of the international terminal. In Concourse E, international passengers ending their journeys in Atlanta will go through a dedicated underground walkway to the Concourse F checkpoint. International passengers arriving in Concourse E that are connecting to another flight will be processed in a separate checkpoint on Concourse E and reenter the concourse via a dedicated TSA checkpoint.
Concourse T contains 21 gates.
Concourse A contains 29 gates.
Concourse B contains 32 gates.
Concourse C contains 34 gates.
Concourse D contains 40 gates.
Concourse E contains 28 gates.
Concourse F contains 12 gates.
Ground transportation
The domestic terminal can be accessed directly from Interstate 85 SB at exit 72/Camp Creek Pkwy, or from Interstate 85 NB at exit 71/Riverdale Rd. The international terminal is accessed directly from Interstate 75 SB or NB at exit 239. These freeways in turn connect with the following additional freeways within 10 miles: Interstate 285, Interstate 675, Georgia State Route 166, Interstate 20.
Hartsfield–Jackson has its own train station on the city's rapid transit system, MARTA, served by the Red and Gold lines. The above-ground station is inside the main building, between the north and south domestic terminals on the west end. The Airport station is currently the southernmost station in the MARTA system, though expansions via metro or commuter rail further south into Clayton County have been discussed.
Several local shared-ride shuttle services are readily available at Atlanta Airport, offering diverse options for travelers seeking convenient transportation.
The Hartsfield–Jackson Rental Car Center, which opened December 8, 2009, houses all ten airport rental agencies with capacity for additional companies. The complex features 9,900 parking spaces split between two four-story parking decks that together cover , a customer service center, and a maintenance center featuring 140 gas pumps and 30 wash bays equipped with a water recovery system. An automated people mover, the ATL SkyTrain, runs between the rental car center, the Domestic Terminal, and the Gateway Center of the Georgia International Convention Center, while a four-lane roadway that spans Interstate 85 connects the rental car center with the existing airport road network.
Other facilities
The 990 Toffie Terrace hangar, a part of Hartsfield–Jackson Airport and located within the City of College Park corporate limits, is owned by the City of Atlanta. The building now houses the Atlanta Police Department Helicopter Unit. It once served as the headquarters of the regional airline ExpressJet.
Before its merger with ExpressJet, Atlantic Southeast Airlines was headquartered in the hangar, then named the A-Tech Center. In December 2007, the airline announced it was moving its headquarters into the facility, previously named the "North Hangar." The hangar includes of hangar bays for aircraft maintenance. It has of adjacent land and 1,400 parking spaces for employees. The airline planned to relocate 100 employees from Macon to the new headquarters. The Atlanta City Council and Mayor of Atlanta Shirley Franklin approved the new 25-year ASA lease, which also gave the airline new hangar space to work on 15 to 25 aircraft in overnight maintenance; previously, its aircraft were serviced at ConcourseC. The airport property division stated that the hangar was built in the 1960s and renovated in the 1970s. Eastern Air Lines and Delta Air Lines had previously occupied the hangar. Delta's lease originally was scheduled to expire in 2010, but the airline returned the lease to the City of Atlanta in 2005 as part of its bankruptcy settlement. The city collected an insurance settlement of almost $900,000 due to the cancellation.
Airlines and destinations
Passenger
: Ethiopian Airlines flights from Addis Ababa to Atlanta stop in Rome–Fiumicino for refueling. The flight from Atlanta to Addis Ababa is nonstop.
Cargo
Statistics
Top destinations
Airline market share
Annual traffic
On-time performance (domestic major U.S. carriers only)
Accidents and incidents
On May 23, 1960, Delta Air Lines Flight 1903, a Convair CV-880-22-1 (N8804E), crashed on takeoff resulting in the loss of all four crew members. This flight was a training flight for two Delta captains who were being type-rated on the 880.
On February 25, 1969, Eastern Air Lines Flight 955 was hijacked by one passenger shortly after takeoff from ATL en route to Miami. The man pulled a .22 caliber pistol and demanded to be flown to Cuba. He got off the plane in Cuba while the DC-8 was allowed to fly back to the U.S.
On April 4, 1977, Southern Airways Flight 242 was on descent to the airport when hail was ingested into the engines, leading them to fail. Pilot errors and difficult weather forced the pilots to attempt an emergency landing on a highway. Upon touchdown, the aircraft struck several buildings and cars, killing 72 people.
On January 18, 1990, Eastern Air Lines Flight 111, a Boeing 727, overran a Beechcraft King Air operated by Epps Air Service, based at another Atlanta airport. The King Air had landed and was taxiing when the 727, still at high speed in its landing roll, collided with the aircraft. The larger plane's wing impacted the roof of the smaller. The pilot of the King Air, an Epps charter pilot, was killed, while a passenger survived. No crew or passengers on the Eastern plane were injured.
On November 1, 1998, AirTran Airways Flight 867, a Boeing 737, lost control and skidded off of the runway while landing, with main landing gear in a drainage ditch and its empennage extending over the taxiway. The nose gear was folded back into the electrical/electronic compartment and turned 90 degrees from its normal, extended position. The cause was an improperly repaired hydraulic line leak that caused the flight crew to lose control of the airplane.
On August 27, 2024, two Delta TechOps employees were killed and another worker seriously injured and taken to hospital when a tire from a Boeing 757 being maintained/disassembled exploded. The tire was not connected to the aircraft at the time. OSHA has opened up an investigation into the incident.
On September 10, 2024, Delta Air Lines Flight 295, an Airbus A350-900, and Endeavor Air Flight 5526, a Bombardier CRJ-900, were involved in a low-speed ground collision. While the CRJ-900 was stationary at a runway hold-short point, the right wingtip of the A350 impacted and broke off the CRJ-900's vertical stabilizer. There were no injuries reported among the 236 people on board the A350 or 59 people on board the CRJ-900.
In popular culture
The airport was seen in the 2018 Channel 5 documentary The Secret Life of the World's Busiest Airport.
| Technology | North America | null |
248457 | https://en.wikipedia.org/wiki/Hong%20Kong%20International%20Airport | Hong Kong International Airport | Hong Kong International Airport is an international airport on the island of Chek Lap Kok in western Hong Kong. The airport is also referred to as Chek Lap Kok International Airport or Chek Lap Kok Airport, to distinguish it from its predecessor, the former Kai Tak Airport.
Opened in 1998, Hong Kong International Airport is the world's busiest cargo gateway and one of the world's busiest passenger airports. It is also home to one of the world's largest passenger terminal buildings, which was the largest when the airport opened.
The airport is operated by Airport Authority Hong Kong (AAHK), a statutory body of the Hong Kong government established on 1 December 1995. It runs 24 hours a day and is the primary hub for Cathay Pacific, Greater Bay Airlines, Hong Kong Airlines, HK Express, and Air Hong Kong (cargo carrier). The airport is one of the hubs of Oneworld, and also one of the Asia-Pacific cargo hubs for UPS Airlines. It is a focus city for China Airlines and China Eastern Airlines. Ethiopian Airlines utilizes Hong Kong as a stopover point for their flights.
Hong Kong International Airport, which employed about 60,000 people at the start of 2024, is an important contributor to Hong Kong's economy. The economic contribution generated by Hong Kong's air travel industry in 2018 amounted to US$33 billion, 10.2% of Hong Kong's GDP. More than 100 airlines operate flights from the airport to over 180 cities across the globe. In 2015, HKIA handled 68.5 million passengers, making it the 8th busiest airport worldwide by passenger traffic and the 4th busiest airport worldwide by international passenger traffic. Since 2010, it has also surpassed Memphis International Airport to become the world's busiest airport by cargo traffic (excluding 2020 due to disruptions related to the COVID-19 pandemic).
History
Chek Lap Kok Airport was designed as a replacement for the former Hong Kong International Airport (commonly known as Kai Tak Airport), built in 1925. Located in the densely built-up Kowloon City District with a single runway extending into Kowloon Bay, Hong Kong Airport had turned on the runway lights for expansion to cope with steadily increasing air traffic. By the 1990s, Kai Tak had become one of the world's busiest airports, being a major hub for multiple passenger airlines along with a major cargo and maintenance hub – it far exceeded its annual passenger and cargo design capacities, and one out of every three flights experienced delays, largely due to a lack of space for aircraft, gates, and a second runway. In addition, noise mitigation measures restricted nighttime flights, as severe noise pollution (exceeding 105 dB(A) in Kowloon City) adversely affected an estimated 340,000 people at least.
A 1974 planning study by the Civil Aviation and Public Works departments identified the small island of Chek Lap Kok, off Lantau Island, as a possible future airport replacement site. Far from the congested city centre, flight paths would be routed over the South China Sea rather than crowded urban areas, enabling efficient round-the-clock operation of multiple runways. The Chek Lap Kok (CLK) airport master plan and civil engineering studies were completed between 1982 and 1983, respectively. However, in February 1983, the government shelved the project for financial and economic reasons. In 1988, the Port & Airport Development Strategy (PADS) study was undertaken by consultants, headed by Mott MacDonald Hong Kong Limited, reporting in December 1989. This study looked at forecasts for both airport and port traffic to the year 2011 and came up with three recommended strategies for overall strategic development in Hong Kong. One of the three assumed maintaining the existing airport at Kai Tak; a second assumed a possible airport in the Western Harbour between Lantau Island and Hong Kong Island, and the third assumed a new airport at Chek Lap Kok. The consultants produced detailed analyses for each scenario, enabling the government to consider these appraisals for each of the three "Recommended Strategies". In October 1989, the Governor of Hong Kong announced to the Legislative Council that a decision had been made on the territory's long-term port and airport development strategy. The strategy was to be adopted that included a replacement airport at Chek Lap Kok and incorporating new container terminals 8 and 9 at Stonecutters Island and east of the Tsing Yi island, respectively.
In the PADS study, the consultants advised that the earliest the airport could be opened was January 1998. However, in reaching the government's decision, this date was changed to January 1997, six months before the handover of Hong Kong to China. Construction of the new airport was began in 1991. As construction progressed, an agreement was reached with China that as much as possible of the airport would be completed before the handover to China in July 1997. British Prime Minister, John Major, opened the Tsing Ma Bridge, the main access to Lantau Island and the airport and its supporting community in May 1997. Soon after, the airport itself opened in July 1998.
The construction period was extremely rushed; specialists considered that only a 10–20 years period was sufficient for this massive project. Another cause for this rush was due to the uncertain future of the airport construction after the transfer of sovereignty over Hong Kong to the People's Republic of China. Shortly after the then-British colonial government of Hong Kong announced plans to construct the new airport, the Chinese government in Beijing began voicing objections to various aspects of the massive project, which prompted financial institutions to delay extending project finance. Without access to this financing, many of the companies who had secured contracts to build various portions of the project halted the construction, resulting in delays that pushed the actual opening of the airport which was originally planned to take place before the transition in sovereignty until one year after. As agreements were reached with the government in China, Beijing withdrew most of its objections and work then continued, albeit behind schedule.
Hong Kong International Airport was built on a large artificial island formed by flattening and levelling the former Chek Lap Kok and Lam Chau islands ( and respectively) and reclaiming of the adjacent seabed. The airport site with its reclamation, added nearly 1% to Hong Kong's total surface area, connecting to the north side of Lantau Island near Tung Chung new town.
Construction of the new airport was only part of the Airport Core Programme, which also involved the construction of new roads and rail links to the airport, with associated bridges and tunnels, and major land reclamation projects on both Hong Kong Island and in Kowloon. The project holds the record for the most expensive airport project ever, according to Guinness World Records. The construction of the new airport was also voted as one of the Top 10 Construction Achievements of the 20th Century at the ConExpo conference in 1999.
The detailed design for the airport terminal was awarded to a consortium led by Mott Connell (the Hong Kong office of UK consultant Mott MacDonald) with British Airports Authority as the specialist designers for airport-related aspects, Foster and Partners as the architects, and Ove Arup as the specialist structural designers for the roof. Mott Connell was the designer for the foundations, all other structural components, and the mechanical and electrical work. The sides of the terminals, predominantly glass, were designed to break during high-speed winds, relieving pressure and allowing the terminal to withstand an intense typhoon.
The airport was officially opened in an opening ceremony by the President of the People's Republic of China and General Secretary of the Communist Party Jiang Zemin at noon Hong Kong Time on 2 July 1998. Hours later, Air Force One, carrying the President of the United States Bill Clinton, landed at the new airport and became the first foreign visitor to arrive at the new airport. The actual operation of the airport commenced on 6 July 1998, concluding the six-year construction that cost 60 billion US dollar. On that day at 06:25 Hong Kong Time, Cathay Pacific Flight CX 889 from New York JFK Airport became the first commercial flight to land at the airport, pipping the original CX 292 from Rome which was the scheduled as the first arrival. However, the airport had already started to experience some technical difficulties on the first day of opening. The flight information display system (FIDS) had suddenly shut down which caused long delays. Shortly afterwards, the cargo-communication link with Kai Tak, where all the necessary data was stored (some still stored there then), went down. During the same period, someone accidentally deleted an important database for cargo services. This meant that cargo had to be manually stored. At one point, the airport had to turn away all air cargo and freight headed for and exported from Hong Kong (except food and medical supplies) while it sorted out the huge mess. HKIA simply could not keep up without an automated assistant-computer system. For three to five months after its opening, it suffered various severe organisational, mechanical, and technical problems that almost crippled the airport and its operations. Computer glitches were mostly to blame for the major crisis. Lau Kong-wah, a Hong Kong politician, was quoted saying, "This was meant to be a first-class project, but it has turned into a ninth-class airport and a disgrace. Our airport has become the laughingstock of the world." At one time, the government reopened the cargo terminal at Kai Tak Airport to handle freight traffic because of a breakdown at the new cargo terminal, named Super Terminal One (ST1). However, after six months, the airport started to operate normally.
On 31 July 2000, Todd Salimuchai, a regularised illegal immigrant in Hong Kong with no provable nationality, forced his way through a security checkpoint using a fake pistol, took a woman hostage, and boarded a Cathay Pacific aircraft. He was demanded to be flown to Burma, which he claimed was his native country but had refused to admit him due to his lack of documents. He surrendered to the police two and a half hours later.
Officially opened in June 2007, the second airport terminal, called T2 (check-in facility only), is linked with the MTR Airport Express on a new platform. The terminal also features a new shopping mall, SkyPlaza, providing a large variety of shops and restaurants, together with a few entertainment facilities. T2 also houses a 36-bay coach station for buses to and from mainland China and 56 airline check-in counters, as well as customs and immigration facilities.
Besides T2, the SkyCity Nine Eagles Golf Course was opened in 2007 whereas the second airport hotel, the Hong Kong SkyCity Marriott Hotel, and a permanent cross-boundary ferry terminal, the Skypier, began operations in 2008 and 2009 respectively. Development around T2 also includes the AsiaWorld-Expo which started operation in late 2005. A second passenger concourse, the North Satellite Concourse (NSC), opened in 2010, followed by the Midfield Concourse in December 2015.
During August 2019, the airport was shut down multiple times as demonstrations were held inside the airport during the 2019–20 Hong Kong protests, over 160 flights were cancelled as both the arrivals and departures sections of the airport were occupied.
The third runway, also known as the North Runway, was opened in July 2022. It is the first part of the Hong Kong International Airport Master Plan 2030 to be implemented. The third runway is 650 hectares in land area, roughly the size of Gibraltar. The Centre Runway and Terminal 2 of the airport were then closed to facilitate construction works, expansion and upgrades. The Center Runway was reopened on the 28th of November 2024. Terminal 2 is expected to re-enter service in 2024.
Composition
Hong Kong International Airport covers an area of 4,707 acres or . The airport has a total of 89 boarding gates, with 77 jet bridge gates (1–22, 23–36, 40–50, 60–71, 201–219) and 12 virtual gates (228–230, 511–513, 520–524) which are used as assembly points for passengers, who are then ferried to the aircraft by apron buses. Of the 77 jet bridges, five (Gates 5, 23, 60, 62, 64) are capable of handling the Airbus A380, the current users of which are Asiana Airlines, British Airways, Emirates, Qantas and Singapore Airlines. Korean Air and China Southern Airlines previously operated a route to HKIA from Seoul and Beijing respectively using the Airbus A380, but these airlines decided to not use them due to unprofitable nature of the aircraft type. Air France, Lufthansa and Thai Airways International previously operated services to Hong Kong from Paris, Frankfurt and Bangkok using the Airbus A380, though they retired the aircraft types early due to the COVID-19 pandemic.
In addition to Chek Lap Kok, the airport occupies what was Lam Chau.
Terminal 1
Terminal 1 of the HKIA, with an area measuring , is one of the largest passenger airport terminal buildings in the world, after the likes of Dubai International Airport Terminal 3 and Beijing Capital International Airport Terminal 3.
Opened on 6 July 1998, Terminal 1 was the largest airport passenger terminal building, with a total gross floor area of . It briefly conceded the status to Bangkok's Suvarnabhumi Airport () when the latter opened on 15 September 2006, but reclaimed the title when the East Hall was expanded, bringing the total area to its current size of . Terminal 1's title as the world's largest was surrendered to Beijing Capital International Airport Terminal 3 on 29 February 2008.
In late 2021, the air side of Terminal 1 started segregating mainland Chinese flights and other international flights into two separate zones, "Green Zone" and "Orange Zone", to reduce the risk of cross infection of novel coronavirus between travellers and airport workers serving different destinations.
On 1 November 2022, the sky bridge opened as part of a wider HK$9 billion airport upgrade, connecting Terminal 1 to the T1 Satellite Concourse (T1S). Lined with glass floor panels at the edges, the 200 metre long and 28 metre high bridge, the largest of its kind, is high enough for an Airbus A380 to pass underneath.
T1 Satellite Concourse
In 2007, HKIA began the construction of a two-story T1 Satellite Concourse (T1S), previously known as the North Satellite Concourse (NSC), which opened in December 2009. This concourse was designed for narrow-body aircraft and is equipped with 10 jet bridges. The concourse has a floor area of and will be able to serve more than five million passengers annually. T1S was built so the airport could accommodate at least 90 percent of its passengers by aerobridges. It has two levels (one for departures and one for arrivals). A new sky bridge connecting Terminal 1 and T1S opened in November 2022, allowing passengers to walk above taxiing planes, saving time from taking the airport shuttle bus.
T1 Midfield Concourse
On 25 January 2011, Airport Authority Hong Kong (AA) unveiled phase 1 of its midfield development project which was targeted for completion by the end of 2015. The midfield area is located to the west of Terminal 1 between the two existing runways. It was the then last piece of land on the airport island available for large-scale development. This includes 20 aircraft parking stands, three of which are wide enough to serve the Airbus A380 and cater for an additional 10 million passengers annually. Passengers reach the concourse through an extension of the underground automated people mover. A joint venture of Mott MacDonald and Arup led the design of the project. Gammon Construction undertook the construction work. The Concourse began operations on 28 December 2015, and the first flight that used it was the HX658 operated by the Hong Kong Airlines flying from Hong Kong to Okinawa. On 31 March 2016, the concourse was officially inaugurated in a ceremony marking its full commissioning.
Former Terminal 2
Former Terminal 2 with an area measuring , together with the SkyPlaza, opened on 28 February 2007 along with the opening of the Airport station's Platform 3. It was only a low-cost carrier check-in and processing facility for departing passengers with no gates or arrival facilities (passengers were transported underground to gates at Terminal 1). The SkyPlaza was situated within. Former Terminal 2 was shut down on 28 November 2019 at 23:00 to make way for a new satellite terminal from the three-runway system.
Other buildings
Cathay Pacific City, the head office of Cathay Pacific and Air Hong Kong, is located on the airport island. CNAC House, the office for Air China is also located in the airport complex, together with the Civil Aviation Department headquarters. HAECO also has its head office on the airport property. HK Express has its head office on the airport property, in what was previously the Dragonair House, head office of Cathay Dragon.
The Government Flying Service (GFS) has its head office building in the airport. The head office of the Air Accident Investigation Authority (AAIA) is in the Facility Building on the airport property.
Airport expansion projects
In June 2010, the Airport Authority unveiled plans to develop in stages the vast midfield site of the airport island. Stage 1 will involve the construction of a new 20-gate passenger concourse to be built in two phases (completion 2015 and 2020) with 11 gates in phase 1 growing to 20 gates in phase 2. The configuration of the new concourse is similar to those at Atlanta, Bangkok–Suvarnabhumi, Berlin (Terminal 1), Chicago–O'Hare (Global Terminal), Denver, Detroit (McNamara Terminal), London–Heathrow (Terminals 2 and 5), Los Angeles (TBIT), Munich (Terminal 2), Salt Lake City, Seoul–Incheon, Washington–Dulles and Felipe Ángeles International Airport (Zumpango). After stage 1 of midfield development is completed in 2020, there will be sufficient lands remaining for further new concourses to be built as and when demand for them materialises.
Master Plan 2030
One year after, on 2 June 2011, the Airport Authority announced and released their latest version of a 20-year blueprint for the airport's development, the Hong Kong International Airport Master Plan 2030. The study took three years and according to the authority, nine consulting organisations have been hired for the research, observation, planning and advice. The main focus is to improve the overall capacity and aircraft handling ability of the airport. Based on this, two options have been developed.
Option 1: Two-runway system
To maintain the two-runway system, there would be enhancements to the terminal and apron facilities to increase the airport's capacity. This option would enable the airport to handle a maximum of 420,000 flight movements per year, with annual passenger and cargo throughput increased to 74 million and six million tonnes respectively. The approximate cost of this plan was $23.4 billion Hong Kong dollars in 2010 prices. The Airport Authority estimated that the airport would reach its maximum runway capacity sometime around 2020 if no extra runway were to be added.
Option 2: Three-runway system
This plan focussed on adding a third runway to the north of the Chek Lap Kok, the island the airport is built on, by land reclamation, using deep cement mixing, of about . Associated facilities, additional terminals, airfield and apron facilities, would be built as well, and, combined with the new runway, it was estimated that the airport would be able to handle a maximum of 620,000 flights per year (102 per hour, or about one flight every 36 seconds), and meet forecast annual passenger and cargo throughput of about 97 million and 8.9 million tonnes by 2030 respectively.
There were possible drawbacks. Development costs were a concern: although the proposal would increase the number of direct jobs associated with HKIA to 150,000 by 2030 and generate an ENPV of HK$912 billion (in 2009 dollars), the estimated cost was approximately $86.2 billion (2010) Hong Kong Dollars. There were also environmental and local noise pollution concerns.
On 20 March 2012, the Hong Kong Government adopted this option as the official expansion plan.
The third runway, with its dedicated passenger concourse (T2 Concourse), was built parallel to the current two runways on reclaimed land directly north of the existing airport island. The third runway (referred to as the North runway) began operations in July 2022, while the original North runway (re-designated as the Centre runway) was closed for reconfiguration until 2024. Other facilities of the Three-runway system project include the T2 expansion, new T2 Concourse, automatic people mover, and baggage handling system.
Airlines and destinations
Passenger
Cargo
Statistics
Operations
The airport is operated by the Airport Authority Hong Kong, a statutory body wholly owned by the Government of Hong Kong Special Administrative Region.
The airport has three parallel runways, all of which are in length and wide. The south runway has a Category II Precision Approach, while the centre runway has the higher Category IIIA rating, which allows pilots to land in only visibility. The two runways have a capacity of over 60 aircraft movements an hour. The airport is upgrading ATC and runways so that they can handle 68 movements per hour. Normally, the centre runway (07C/25C; until 1 December 2021 the north runway 07L/25R) is used for landing passenger planes. The south runway (07R/25L) is used for passenger planes taking off and cargo flights due to its proximity to the cargo terminal. A third runway (designated 07L/25R) to their north has opened in mid-2022, while the Centre runway has been closed for upgrades. When all three runways are opened, it is estimated that the airport will be able to handle a maximum of 620,000 flights per year (102 per hour, or about one flight every 36 seconds).
There are 49 frontal stands at the main passenger concourse, 28 remote stands and 25 cargo stands. There are also five parking bays at the Northwest Concourse. A satellite concourse with 10 frontal stands for narrow-body aircraft was commissioned to the north of the main concourse at the end of 2009, bringing the total number of frontal stands at the airport to 59.
The airport was the busiest for passenger traffic in Asia in 2010, and the world's busiest airport for cargo traffic in 2010. In terms of international traffic, the airport is the third busiest for passenger traffic and the busiest for cargo since its operation in 1998. Over 95 international airlines are providing about 900 scheduled passenger and all-cargo flights each day between Hong Kong and some 160 destinations worldwide. About 76 percent of these flights are operated with wide-body jets. There is also an average of approximately 31 non-scheduled passenger and cargo flights each week.
The operation of scheduled air services to and from Hong Kong is facilitated by air services agreements between Hong Kong and other countries. Since the opening of HKIA, the Government of the Hong Kong Special Administrative Region has implemented a policy of progressive liberalisation of air services. Many low-cost airlines have started various regional routes to compete head-on with full-service carriers on trunk routes.
The airport's long-term expansion opportunities are subject to variables. The airport opened its third runway in July 2022 as part of a HK$141.5 billion expansion project that would increase its land footprint by 50%. On the other hand, there exists only one airway between Hong Kong and mainland China, and this single route is often and easily backed up causing delays on both sides. In addition, China requires that aircraft flying the single air route between Hong Kong and the mainland must be at an altitude of at least 15,000 feet. Talks are underway to persuade the Chinese military to relax its airspace restriction because of worsening air traffic congestion at the airport. Other than that, Hong Kong Airport Authority is cooperating with other airports in the area to relieve air traffic and in the future, Shenzhen may act as a regional airport while Hong Kong receives all the international flights.
Air traffic
The Government Flying Service provides short and long-range search and rescue services, police support, medical evacuation and general-purpose flights for the Government.
Passenger facilities
Despite its size, the passenger terminal was designed for convenience. The layout and signage, moving walkways and the automated people mover help passengers move through the building. The HKIA Automated People Mover, a driverless people mover system with 3 stations transports passengers between the check-in area and the gates. The trains travel at . The airport also contains an IMAX theatre that has the largest screen in Hong Kong. The theatre is located in Terminal 2, level 6 and can seat 350 persons at a time.
Hong Kong Business Aviation Centre
The Hong Kong Business Aviation Centre (HKBAC) is located within the airport and has its terminal and facilities separate from the public terminal. It provides services for executive aircraft and passengers, including a passenger lounge, private rooms and showers, business centre facilities, ground handling, baggage handling, fuelling, security, customs and flight planning. Designated spaces and hangars are also provided at the HKBAC for private aircraft. HKBAC has broken ground on a HK$400 million ($51 million) expansion. The project, which will double the airport's handling capacity for business jet movements, is expected to be completed in 2025.
Intermodal transportation hub
To sustain the growth of passengers, the Airport Authority formulated a "push and pull through" strategy to expand its connections to new sources of passengers and cargo. This means adapting the network to the rapidly growing markets in China and in particular to the Pearl River Delta region (PRD). In 2003, a new Airport-Mainland Coach Station opened. The coach station has a waiting lounge and sheltered bays for ten coaches. Many buses operate each day to transport passengers between HKIA and major cities on the Mainland.
The Coach Station was relocated to the ground floor (level 3) of Terminal 2 in 2007. The 36 bays at the new Coach Station allow cross-border coaches to make 320 trips a day carrying passengers between the airport and 90 cities and towns in the PRD. Local tour and hotel coaches also operate from T2. The coach station at T2 has shops and waiting lounges as well as a mainland coach service centre which gathers all operators together.
In late September 2003, the SkyPier high-speed ferry terminal opened. Passengers arriving at the SkyPier board buses to the terminal and arriving air passengers board ferries at the pier for their ride back to the PRD. Passengers travelling in both directions can bypass customs and immigration formalities, which reduces transit time. Four ports – Shekou, Shenzhen, Macau and Humen (Dongguan) – were initially served. As of August 2007, SkyPier serves Shenzhen's Shekou and Fuyong, Dongguan's Humen, Macau, Zhongshan and Zhuhai. Passengers travelling from Shekou and Macau can complete airline check-in procedures with participating airlines before boarding the ferries and go straight to the boarding gate for the flight at HKIA.
In 2009, the permanent SkyPier Terminal opened. The permanent ferry terminal is equipped with four berths, but the terminal is designed to accommodate eight berths. Transfer desks and baggage handling facilities are included, and the terminal is directly connected to the airport automatic people mover system.
Baggage and cargo facilities
Ramp handling services are provided by Hong Kong Airport Services Limited (HAS), Jardine Air Terminal Services Limited and SATS HK Limited. Their services include the handling of mail and passenger baggage, transportation of cargo, aerobridge operations and the operation of passenger stairways. The airport has an advanced baggage handling system (BHS), the main section of which is located in the basement level of the passenger terminal, and a separate remote transfer facility at the western end of the main concourse for the handling of tight connection transfer bags.
HKIA handles over five million tonnes of cargo annually. Hong Kong Air Cargo Terminals Limited operates one of the two air cargo terminals at the airport. Its headquarters, the SuperTerminal 1, is the world's second-largest stand-alone air cargo handling facility, after the opening of the West Cargo Handling Area of the Shanghai Pudong International Airport on 26 March 2008. The designed capacity is 2.6 million tonnes of freight a year. The second air cargo terminal is operated by Asia Airfreight Terminal Company Limited, and has a capacity of 1.5 million tonnes a year. DHL operates the DHL Central Asia Hub cargo facility which handles 35,000 parcels and 40,000 packages per hour. Hongkong Post operates the Air Mail Centre (AMC) and processes 700,000 packages per day. It is envisaged that HKIA's total air cargo capacity per annum will reach nine million tonnes ultimately.
Aircraft maintenance services
Both line and base maintenance services are undertaken by Hong Kong Aircraft Engineering Company (HAECO), while China Aircraft Services Limited (CASL) and Pan Asia Pacific Aviation Services Limited carry out line maintenance. Line maintenance services include routine servicing of aircraft performed during normal turnaround periods and regularly scheduled layover periods. Base maintenance covers all airframe maintenance services and for this HAECO has a three-bay hangar, which can accommodate up to three Boeing 747-400 aircraft and two Airbus A320 aircraft, and an adjoining support workshop. HAECO also has the world's largest mobile hangar, weighing over 400 tons. It can be used to enclose half of a wide-body aeroplane so that the whole facility can fully enclose four 747s when the mobile hangar is used.
On 29 May 2009, CASL opened its first aircraft maintenance hangar in the maintenance area of the airport. The new hangar occupies an area of about and can accommodate one wide-body and one narrow-body aircraft at the same time; the hangar also has an about area in its annexe building. CASL specialises in Airbus A320 family and Boeing 737 Next Generation series heavy maintenance.
Airport based ground services
The Air Traffic Control Complex (ATCX), located at the centre of the airfield, is the nerve centre of the entire air traffic control system. Some 370 air traffic controllers and supporting staff work around the clock to provide air traffic control services for the Hong Kong Flight Information Region (FIR). At the Air Traffic Control Tower, controllers provide 24-hour aerodrome control services to aircraft operating at the airport. A backup Air Traffic Control Centre/Tower constructed to the north of the ATCX is available for operational use in the event normal services provided in the ATCX are disrupted by unforeseen circumstances. Apart from serving as an operational backup, the facilities are also used for air traffic control training.
The Airport Meteorological Office (AMO) of the Hong Kong Observatory (HKO) provides weather services for the aviation community. The AMO issues alerts of low-level windshear and turbulence. Windshear detection is made using traditional doppler weather radars as well as the more effective doppler LIDAR, of which Hong Kong International Airport was the first to introduce. Doppler LIDAR systems use lasers to detect windshear and wind direction even when atmospheric conditions are too dry for Doppler radar to work.
Fire and rescue services
Rescue and fire fighting services within the airport are covered by the Airport Fire Contingent of the Hong Kong Fire Services Department. The contingent has 282 members, operating three fire stations and two rescue berths for 24-hour emergency calls. It is equipped with 14 fire appliances which can respond to incidents within two minutes in optimum conditions of visibility and surface conditions, satisfying the relevant recommendation of the International Civil Aviation Organization. Two high-capacity rescue boats, supported by eight-speed boats, form the core of sea rescue operations. One ambulance is assigned at each of the airport fire stations.
Ground transport
The Airport is connected to inner Hong Kong by the Route 8 in Hong Kong North Lantau Highway on Lantau Island.
There is an automated people mover, operated by the Airport Authority and maintained by MTR Corporation, connecting the East Hall to the Midfield Concourse via West Hall and Terminal 2. It was extended to SkyPier in late 2009 and extended to Midfield Concourse in 2015.
Bus
Citybus (CityFlyer for Airport services), New Lantau Bus, Long Win Bus and Discovery Bay Transit Services (Permits required) operate more than 40 bus routes to the airport from various parts of Hong Kong, available at the Airport Ground Transportation Centre and Cheong Tat Road. The bus companies also offer more than 20 overnight "N" and "NA" Bus lines (a.k.a. Night services).
Passengers can also take bus route number S1 to the Tung Chung MTR station. From there they can board the MTR Tung Chung line which follows the same route as the MTR Airport Express Line to Central Station with cheaper fare but longer journey time.
There is a bus service to Hong Kong–Zhuhai–Macau Bridge Control Point, with services between Chek Lap Kok, Hong Kong to Zhuhai and Macau. Coach services are also available to major cities and towns in Guangdong province. such as Dongguan, Guangzhou and Shenzhen. And Also for HZMBus to Macau
Ferry
Direct ferry services are available from the airport to various destinations throughout the Pearl River Delta (including Macau) via Skypier. Passengers using these services are treated as transit passengers and are not considered to have entered Hong Kong for immigration purposes. For this reason, access to the ferry terminal is before immigration at the airport for arriving passengers. Check-in services are available at these piers. Four ports – Shekou, Shenzhen Airport (Fuyong) and Humen (Dongguan) in mainland China, and Outer Harbour Ferry Terminal in Macau– were initially served, extending to Guangzhou and Zhongshan at the end of 2003. The Zhuhai service began on 10 July 2007 while a Nansha service started on 14 July 2009.
Rail
The fastest service from the city to the airport is the Airport Express, which was a part of the Hong Kong Rail Network in Hong Kong. A dedicated rapid-speed rail link as part of the MTR rapid transit network. The line serves between Asia-World-Expo and Hong Kong (Central) Station makes intermediate stops at the following stations:
Tsing Yi Station (Located in the northeastern part of Tsing Yi Island, Kwai Tsing District, Tsing Yi [which is one of three Communes which form the Tsuen Wan New Town {the other two are Tsuen Wan and Kwai Chung}], this station mostly served passengers from the western part of the New Territories. Transfers are available for the Tung Chung Line. Connections are also available for taxis and public/private buses at the local Maritime Square)
Kowloon Station (Located in the Yau Tsim Mong District on the western part of the Kowloon Peninsula, this station is the major transfer hub in the Kowloon Peninsula, with stunning landmarks such as West Kowloon Cultural District, the M+ Art Museum, Hong Kong Palace Museum, Avenue of Stars and many more. Transfers are available for the Tung Chung Line [and Tuen Ma Line (Formerly called the West Rail Line) in Austin Station] for passengers for the East Kowloon and the New Territories. Since Autumn 2018, the High Speed Rail Network Line operates in West Kowloon Station which connects to the National Rail Network of China. Connections are also available for taxis, MTR Shuttle Buses and public/private buses at Elements. Before the COVID-19 pandemic and currently suspended under further notice as of April 2023, in-town check-in services for major airlines were provided.
Hong Kong Station, the terminus, is located at the northern coast of Central and Western District on Hong Kong Island. It takes approximately 24 minutes to reach the airport from Hong Kong Station. Transfers are available for the Tung Chung Line (and Island Line and Tsuen Wan Line at Central Station, where travellators link the two stations.) Connections are also available such as free MTR shuttle bus services between Airport Express stations and hotels in the area, and free transfers to and from other MTR lines with a valid Octopus card (which is not available to Single Ride Ticket users). Hong Kong Station also provides in-town check-in services for major airlines. (Passengers can ride one stop of the Island Line or Tsuen Wan Line to Admiralty Station, where transfers are available for the South Island Line [Opened on December 28, 2016, located on Platforms 5 and 6 {towards South Horizons}] and the East Rail Line [Opened on May 15, 2022, located on Platforms 7 {towards Lo Wu and Lok Ma Chau Border Crossing Stations} and 8 {terminating platform}.])
The Airport Express line originally terminated at Airport station, where trains open doors on both sides, allowing direct access to either Terminal 1 or Terminal 2. It was later extended to AsiaWorld–Expo station on 20 December 2005 to facilitate the opening of the nearby AsiaWorld–Expo venue. During events at the venue, some Tung Chung line trains, which largely share the same tracks as the Airport Express, serve this station instead of Tung Chung, but these trains do not stop by the Airport station.
Taxi
The airport is served by three types of taxis, distinguished by colours:
Urban Taxis connect the Airport with Hong Kong Island, Kowloon Peninsula and parts of the new towns of Metropolitan Hong Kong such as Tsuen Wan, Sha Tin and Tseung Kwan O.(urban taxis can go anywhere in Hong Kong except southern parts of Lantau Island).
New Territories Taxis connect the airport with the New Territories, except those parts in the Metropolitan Hong Kong Area such as Tsuen Wan, Sha Tin and Tseung Kwan O (except parts of Hang Hau) were served by urban taxis.
Lantau Taxis connect the airport with the rest of Lantau Island.
Accidents and incidents
The following are aviation accidents or incidents at the current HKIA (see accidents and incidents at the former HKIA at Kai Tak):
On 22 August 1999, China Airlines Flight 642 (an MD-11 operated by subsidiary Mandarin Airlines), which was landing at Hong Kong International Airport during Typhoon Sam after a flight from Bangkok International Airport (now Bangkok Don Mueang International Airport), rolled over and caught fire, coming to rest upside down beside the runway. Of the 315 passengers and crew on board, 3 people were killed and 219 were injured.
On 13 April 2010, Cathay Pacific Flight 780, an Airbus A330-342 from Surabaya Juanda International Airport to Hong Kong landed safely after both engines failed due to contaminated fuel. All 322 survived, and 63 of them were injured. Its two pilots received the Polaris Award from the International Federation of Air Line Pilots' Associations for their heroism and airmanship.
On 8 September 2016, an airport delivery van crashed into the left engine of KA691 from Hong Kong to Penang with 295 passengers and crew on board as the aircraft was taxiing to the runway. There were no fatalities.
On 4 October 2017, some cargo caught fire as it was being loaded onto an American Airlines plane. The authorities were able to contain the fire. The plane and an airport vehicle was damaged and one person was injured.
In April 2021, three cargo pallets waiting to be transferred onto a Hong Kong Air Cargo flight caught fire. Each pallet included Vivo smartphones destined for Bangkok. The airline banned shipments of Vivo phones and all shipments from two freight forwarding companies as a result.
On 21 July 2021, UPS Airlines Flight 5X003, a Boeing 747-8F, suffered from an engine fire in the #1 engine shortly after take-off. The flight subsequently made a successful emergency landing on Runway 07L and was met by airport fire crews from the Hong Kong Fire Services Department Main Airport Station which extinguished the flames. The Hong Kong Civil Aviation Department and the NTSB are currently investigating the incident.
On 20 December 2022, United Airlines Flight 2831 suffered a bird strike on takeoff which triggered a fire. The plane was ordered back and landed on Runway 07L and was met by airport fire crews who extinguished the flames. The Hong Kong Civil Aviation Department and the NTSB are currently investigating the incident.
On 17 June 2024, Atlas Air Flight 4304 suffered a tire burst while performing an emergency landing. The tire fragments caused a 4-hour delay that delayed 186 of the 315 flights scheduled during this period. The Airport Authority reported that none of the five crew members on board the cargo plane suffered injuries. The break down of the hydraulic system caused clean-up to take longer than expected.
Accolades
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248669 | https://en.wikipedia.org/wiki/Tapir | Tapir | Tapirs ( ) are large, herbivorous mammals belonging to the family Tapiridae. They are similar in shape to a pig, with a short, prehensile nose trunk. Tapirs inhabit jungle and forest regions of South and Central America and Southeast Asia. They are one of three extant branches of Perissodactyla (odd-toed ungulates), alongside equines and rhinoceroses. Only a single genus, Tapirus, is currently extant. Tapirs migrated into South America during the Pleistocene epoch from North America after the formation of the Isthmus of Panama as part of the Great American Interchange. Tapirs were formerly present across North America, but became extinct in the region at the end of the Late Pleistocene, around 12,000 years ago.
Name
The term tapir comes from the Portuguese-language words , , which themselves trace their origins back to Old Tupi, specifically the term . This word, according to Eduardo de Almeida Navarro, referred in a more precise manner to the species Tapirus terrestris.
Species
There are four widely recognized extant species of tapir, all in the genus Tapirus of the family Tapiridae. They are the South American tapir, the Malayan tapir, Baird's tapir, and the mountain tapir. In 2013, a group of researchers said they had identified a fifth species of tapir, the kabomani tapir. However, the existence of the kabomani tapir as a distinct species has been widely disputed, and recent genetic evidence further suggests that it actually is part of the species South American tapir. Extant species
The four species are all classified on the IUCN Red List as Endangered or Vulnerable. The tapirs have a number of extinct relatives in the superfamily Tapiroidea. The closest extant relatives of the tapirs are the other odd-toed ungulates, which include horses, wild asses, zebras and rhinoceroses.
Extinct species
During the Late Pleistocene, several other species inhabited North America, including Tapirus veroensis, native to the southern and eastern United States (with its northernmost records being New York State), and Tapirus merriami and Tapirus californicus, native to Western North America. These became extinct during the Quaternary extinction event around 12,000 years ago, along with most of the other large mammals of the Americas, co-inciding with the first arrival of humans to the continent. Tapirus augustus (formerly placed in Megatapirus), native to Southeast and East Asia, substantially larger than the Malayan tapir, also became extinct at some point during the Late Pleistocene. Many primitive tapirs were originally classified under Palaeotapirus including members of Paratapirus and Plesiotapirus, but the original diagnostic material of the genus was too poor to characterize, leading to included species being moved to new genera.
General appearance
Size varies between types, but most tapirs are about long, stand about high at the shoulder, and weigh between . Their coats are short and range in colour from reddish brown, to grey, to nearly black, with the notable exceptions of the Malayan tapir, which has a white, saddle-shaped marking on its back, and the mountain tapir, which has longer, woolly fur. All tapirs have oval, white-tipped ears, rounded, protruding rumps with stubby tails, and splayed, hooved toes, with four toes on the front feet and three on the hind feet, which help them to walk on muddy and soft ground. Baby tapirs of all types have striped-and-spotted coats for camouflage.
Females have a single pair of mammary glands, and males have long penises relative to their body size.
Physical characteristics
The proboscis of the tapir is a highly flexible organ, able to move in all directions, allowing the animals to grab foliage that would otherwise be out of reach. Tapirs often exhibit the flehmen response, a posture in which they raise their snouts and show their teeth to detect scents. This response is frequently exhibited by bulls sniffing for signs of other males or females in oestrus in the area. The length of the proboscis varies among species; Malayan tapirs have the longest snouts and Brazilian tapirs have the shortest. The evolution of tapir probosces, made up almost entirely of soft tissues rather than bony internal structures, gives the Tapiridae skull a unique form in comparison to other perissodactyls, with a larger sagittal crest, orbits positioned more rostrally, a posteriorly telescoped cranium, and a more elongated and retracted nasoincisive incisure.Colbert, Matthew (2002) Tapirus terrestris. Digital Morphology. Retrieved June 20, 2006.
Tapirs have brachyodont, or low-crowned teeth, that lack cementum. Their dental formula is:
Totaling 42 to 44 teeth, this dentition is closer to that of equids, which may differ by one less canine, than their other perissodactyl relatives, rhinoceroses.Huffman, Brent. Order Perissodactyla at Ultimate Ungulate Their incisors are chisel-shaped, with the third large, conical upper incisor separated by a short gap from the considerably smaller canine. A much longer gap is found between the canines and premolars, the first of which may be absent. Tapirs are lophodonts, and their cheek teeth have distinct lophs (ridges) between protocones, paracones, metacones and hypocones.Myers, P., R. Espinosa, C. S. Parr, T. Jones, G. S. Hammond, and T. A. Dewey. 2006. The Basic Structure of Cheek Teeth. The Animal Diversity Web (online). Retrieved June 20, 2006.
Tapirs have brown eyes, often with a bluish cast to them, which has been identified as corneal cloudiness, a condition most commonly found in Malayan tapirs. The exact etiology is unknown, but the cloudiness may be caused by excessive exposure to light or by trauma.Janssen, Donald L., DVM, Dipl ACZM, Bruce A. Rideout, DVM, PhD, Dipl ACVP, Mark E. Edwards, PhD. "Medical Management of Captive Tapirs (Tapirus sp.)." 1996 American Association of Zoo Veterinarians Proceedings. Nov 1996. Puerto Vallarta, Mexico. Pp. 1–11 However, the tapir's sensitive ears and strong sense of smell help to compensate for deficiencies in vision.
Tapirs have simple stomachs and are hindgut fermenters that ferment digested food in a large cecum.
Life cycle
Young tapirs reach sexual maturity between three and five years of age, with females maturing earlier than males. Under good conditions, a healthy female tapir can reproduce every two years; a single young, called a calf, is born after a gestation of about 13 months. The natural lifespan of a tapir is about 25 to 30 years, both in the wild and in zoos. Apart from mothers and their young offspring, tapirs lead almost exclusively solitary lives.
Behaviour
Although they frequently live in dryland forests, tapirs with access to rivers spend a good deal of time in and under water, feeding on soft vegetation, taking refuge from predators, and cooling off during hot periods. Tapirs near a water source will swim, sink to the bottom, and walk along the riverbed to feed, and have been known to submerge themselves to allow small fish to pick parasites off their bulky bodies. Along with freshwater lounging, tapirs often wallow in mud pits, which helps to keep them cool and free of insects.
In the wild, the tapir's diet consists of fruit, berries, and leaves, particularly young, tender vegetation. Tapirs will spend many of their waking hours foraging along well-worn trails, snouts to the ground in search of food. Baird's tapirs have been observed to eat around 40 kg (85 lb) of vegetation in one day.
Tapirs are largely nocturnal and crepuscular, although the smaller mountain tapir of the Andes is generally more active during the day than its congeners. They have monocular vision.
Copulation may occur in or out of water. In captivity, mating pairs will often copulate several times during oestrus. Intromission lasts between 10 and 20 minutes.
Habitat, predation, and vulnerability
Adult tapirs are large enough to have few natural predators, and the thick skin on the backs of their necks helps to protect them from threats such as jaguars, crocodiles, anacondas, and tigers. The creatures are also able to run fairly quickly, considering their size and cumbersome appearance, finding shelter in the thick undergrowth of the forest or in water. Hunting for meat and hides has substantially reduced their numbers and, more recently, habitat loss has resulted in the conservation watch-listing of all four species; the Brazilian tapir is classified as vulnerable, and Baird's tapir, the mountain tapir, and the Malayan tapir are endangered. According to 2022 study published in the Neotropical Biology and Conservation, the lowland tapir in the Atlantic Forest is at risk of complete extinction as a result of anthropogenic pressures, in particular hunting, deforestation and population isolation.
Evolution and natural history
Tapirs originated from the "tapiroids", a group of primitive perissodactyls that inhabited North America and Asia during the Eocene epoch, with tapirs probably originating from the family Helaletidae. The oldest known members of the family Tapiridae such as Protapirus are known from the Early Oligocene of Europe. The oldest representatives of the modern genus Tapirus appeared in Europe during the Mid-Miocene, with Tapirus dispersing into Asia and North America by the late Miocene. Tapirus became extinct in Europe around the end of the Pliocene. Tapirs dispersed into South America during Pleistocene as part of the Great American Biotic Interchange with their oldest records on the continent dating to around 2.6-1 million years ago.
Approximate divergence times based on a 2013 analysis of mtDNA sequences are 0.5 Ma for T. kabomani and the T. terrestris–T. pinchaque clade, 5 Ma for T. bairdii and the three South American tapirs, and 9 Ma for the branching of T. indicus. T. pinchaque arises from within a paraphyletic complex of T. terrestris'' populations.
Genetics
The species of tapir have the following chromosomal numbers:
The Malayan tapir, the species most isolated geographically and genetically, has a significantly smaller number of chromosomes and has been found to share fewer homologies with the three types of American tapirs. A number of conserved autosomes (13 between karyotypes of Baird's tapir and the South American tapir, and 15 between Baird's and the mountain tapir) have also been found in the American species that are not found in the Asian animal. However, geographic proximity is not an absolute predictor of genetic similarity; for instance, G-banded preparations have revealed Malayan, Baird's and South American tapirs have identical X chromosomes, while mountain tapirs are separated by a heterochromatic addition/deletion.
Lack of genetic diversity in tapir populations has become a major source of concern for conservationists. Habitat loss has isolated already small populations of wild tapirs, putting each group in greater danger of dying out completely. Even in zoos, genetic diversity is limited; all captive mountain tapirs, for example, are descended from only two founder individuals.
Hybrids of Baird's and the South American tapirs were bred at the San Francisco Zoo around 1969 and later produced a backcross second generation.
Conservation
A number of conservation projects have been started around the world. The Tapir Specialist Group, a unit of the IUCN Species Survival Commission, strives to conserve biological diversity by stimulating, developing, and conducting practical programs to study, save, restore, and manage the four species of tapir and their remaining habitats in Central and South America and Southeast Asia.
The Baird's Tapir Project of Costa Rica, begun in 1994, is the longest ongoing tapir project in the world. It involves placing radio collars on tapirs in Costa Rica's Corcovado National Park to study their social systems and habitat preferences.
The Lowland Lowland Tapir Conservation Initiative is a conservation and research organization founded by Patrícia Medici, focused on tapir conservation in Brazil.
Attacks on humans
Tapirs are generally shy, but when scared they can defend themselves with their very powerful jaws. In 1998, a zookeeper in Oklahoma City was mauled and had an arm severed after opening the door to a female tapir's enclosure to push food inside (the tapir's two-month-old baby also occupied the cage at the time). In 2006, Carlos Manuel Rodriguez Echandi (who was then the Costa Rican Environmental Minister) became lost in the Corcovado National Park and was found by a search party with a "nasty bite" from a wild tapir. In 2013, a two-year-old girl suffered stomach and arm injuries after being mauled by a South American tapir in Dublin Zoo during a supervised experience in the tapir enclosure. Dublin Zoo pleaded guilty to breaching health and safety regulations and was ordered to pay €5,000 to charity. However, such examples are rare; for the most part, tapirs are likely to avoid confrontation in favour of running from predators, hiding, or, if possible, submerging themselves in nearby water until a threat is gone.
Frank Buck wrote about an attack by a tapir in 1926, which he described in his book, Bring 'Em Back Alive.
Folklore
Tapirs feature in the folklore of several cultures around the world. In Japan, tapirs are associated with the mythological Baku, believed to ward off nightmares. In South America, tapirs are associated with the creation of the earth.
| Biology and health sciences | Perissodactyla | null |
248671 | https://en.wikipedia.org/wiki/Gluconeogenesis | Gluconeogenesis | Gluconeogenesis (GNG) is a metabolic pathway that results in the biosynthesis of glucose from certain non-carbohydrate carbon substrates. It is a ubiquitous process, present in plants, animals, fungi, bacteria, and other microorganisms. In vertebrates, gluconeogenesis occurs mainly in the liver and, to a lesser extent, in the cortex of the kidneys. It is one of two primary mechanisms – the other being degradation of glycogen (glycogenolysis) – used by humans and many other animals to maintain blood sugar levels, avoiding low levels (hypoglycemia). In ruminants, because dietary carbohydrates tend to be metabolized by rumen organisms, gluconeogenesis occurs regardless of fasting, low-carbohydrate diets, exercise, etc. In many other animals, the process occurs during periods of fasting, starvation, low-carbohydrate diets, or intense exercise.
In humans, substrates for gluconeogenesis may come from any non-carbohydrate sources that can be converted to pyruvate or intermediates of glycolysis (see figure). For the breakdown of proteins, these substrates include glucogenic amino acids (although not ketogenic amino acids); from breakdown of lipids (such as triglycerides), they include glycerol, odd-chain fatty acids (although not even-chain fatty acids, see below); and from other parts of metabolism that includes lactate from the Cori cycle. Under conditions of prolonged fasting, acetone derived from ketone bodies can also serve as a substrate, providing a pathway from fatty acids to glucose. Although most gluconeogenesis occurs in the liver, the relative contribution of gluconeogenesis by the kidney is increased in diabetes and prolonged fasting.
The gluconeogenesis pathway is highly endergonic until it is coupled to the hydrolysis of ATP or GTP, effectively making the process exergonic. For example, the pathway leading from pyruvate to glucose-6-phosphate requires 4 molecules of ATP and 2 molecules of GTP to proceed spontaneously. These ATPs are supplied from fatty acid catabolism via beta oxidation.
Precursors
In humans the main gluconeogenic precursors are lactate, glycerol (which is a part of the triglyceride molecule), alanine and glutamine. Altogether, they account for over 90% of the overall gluconeogenesis.
Other glucogenic amino acids and all citric acid cycle intermediates (through conversion to oxaloacetate) can also function as substrates for gluconeogenesis. Generally, human consumption of gluconeogenic substrates in food does not result in increased gluconeogenesis.
In ruminants, propionate is the principal gluconeogenic substrate. In nonruminants, including human beings, propionate arises from the β-oxidation of odd-chain and branched-chain fatty acids, and is a (relatively minor) substrate for gluconeogenesis.
Lactate is transported back to the liver where it is converted into pyruvate by the Cori cycle using the enzyme lactate dehydrogenase. Pyruvate, the first designated substrate of the gluconeogenic pathway, can then be used to generate glucose. Transamination or deamination of amino acids facilitates entering of their carbon skeleton into the cycle directly (as pyruvate or oxaloacetate), or indirectly via the citric acid cycle. The contribution of Cori cycle lactate to overall glucose production increases with fasting duration. Specifically, after 12, 20, and 40 hours of fasting by human volunteers, the contribution of Cori cycle lactate to gluconeogenesis was 41%, 71%, and 92%, respectively.
Whether even-chain fatty acids can be converted into glucose in animals has been a longstanding question in biochemistry. Odd-chain fatty acids can be oxidized to yield acetyl-CoA and propionyl-CoA, the latter serving as a precursor to succinyl-CoA, which can be converted to oxaloacetate and enter into gluconeogenesis. In contrast, even-chain fatty acids are oxidized to yield only acetyl-CoA, whose entry into gluconeogenesis requires the presence of a glyoxylate cycle (also known as glyoxylate shunt) to produce four-carbon dicarboxylic acid precursors. The glyoxylate shunt comprises two enzymes, malate synthase and isocitrate lyase, and is present in fungi, plants, and bacteria. Despite some reports of glyoxylate shunt enzymatic activities detected in animal tissues, genes encoding both enzymatic functions have only been found in nematodes, in which they exist as a single bi-functional enzyme. Genes coding for malate synthase alone (but not isocitrate lyase) have been identified in other animals including arthropods, echinoderms, and even some vertebrates. Mammals found to possess the malate synthase gene include monotremes (platypus) and marsupials (opossum), but not placental mammals.
The existence of the glyoxylate cycle in humans has not been established, and it is widely held that fatty acids cannot be converted to glucose in humans directly. Carbon-14 has been shown to end up in glucose when it is supplied in fatty acids, but this can be expected from the incorporation of labelled atoms derived from acetyl-CoA into citric acid cycle intermediates which are interchangeable with those derived from other physiological sources, such as glucogenic amino acids. In the absence of other glucogenic sources, the 2-carbon acetyl-CoA derived from the oxidation of fatty acids cannot produce a net yield of glucose via the citric acid cycle, since an equivalent two carbon atoms are released as carbon dioxide during the cycle. During ketosis, however, acetyl-CoA from fatty acids yields ketone bodies, including acetone, and up to ~60% of acetone may be oxidized in the liver to the pyruvate precursors acetol and methylglyoxal. Thus ketone bodies derived from fatty acids could account for up to 11% of gluconeogenesis during starvation. Catabolism of fatty acids also produces energy in the form of ATP that is necessary for the gluconeogenesis pathway.
Location
In mammals, gluconeogenesis has been believed to be restricted to the liver, the kidney, the intestine, and muscle, but recent evidence indicates gluconeogenesis occurring in astrocytes of the brain. These organs use somewhat different gluconeogenic precursors. The liver preferentially uses lactate, glycerol, and glucogenic amino acids (especially alanine) while the kidney preferentially uses lactate, glutamine and glycerol. Lactate from the Cori cycle is quantitatively the largest source of substrate for gluconeogenesis, especially for the kidney. The liver uses both glycogenolysis and gluconeogenesis to produce glucose, whereas the kidney only uses gluconeogenesis. After a meal, the liver shifts to glycogen synthesis, whereas the kidney increases gluconeogenesis. The intestine uses mostly glutamine and glycerol.
Propionate is the principal substrate for gluconeogenesis in the ruminant liver, and the ruminant liver may make increased use of gluconeogenic amino acids (e.g., alanine) when glucose demand is increased. The capacity of liver cells to use lactate for gluconeogenesis declines from the preruminant stage to the ruminant stage in calves and lambs. In sheep kidney tissue, very high rates of gluconeogenesis from propionate have been observed.
In all species, the formation of oxaloacetate from pyruvate and TCA cycle intermediates is restricted to the mitochondrion, and the enzymes that convert Phosphoenolpyruvic acid (PEP) to glucose-6-phosphate are found in the cytosol. The location of the enzyme that links these two parts of gluconeogenesis by converting oxaloacetate to PEP – PEP carboxykinase (PEPCK) – is variable by species: it can be found entirely within the mitochondria, entirely within the cytosol, or dispersed evenly between the two, as it is in humans. Transport of PEP across the mitochondrial membrane is accomplished by dedicated transport proteins; however no such proteins exist for oxaloacetate. Therefore, in species that lack intra-mitochondrial PEPCK, oxaloacetate must be converted into malate or aspartate, exported from the mitochondrion, and converted back into oxaloacetate in order to allow gluconeogenesis to continue.
Pathway
Gluconeogenesis is a pathway consisting of a series of eleven enzyme-catalyzed reactions. The pathway will begin in either the liver or kidney, in the mitochondria or cytoplasm of those cells, this being dependent on the substrate being used. Many of the reactions are the reverse of steps found in glycolysis.
Gluconeogenesis begins in the mitochondria with the formation of oxaloacetate by the carboxylation of pyruvate. This reaction also requires one molecule of ATP, and is catalyzed by pyruvate carboxylase. This enzyme is stimulated by high levels of acetyl-CoA (produced in β-oxidation in the liver) and inhibited by high levels of ADP and glucose.
Oxaloacetate is reduced to malate using NADH, a step required for its transportation out of the mitochondria.
Malate is oxidized to oxaloacetate using NAD+ in the cytosol, where the remaining steps of gluconeogenesis take place.
Oxaloacetate is decarboxylated and then phosphorylated to form phosphoenolpyruvate using the enzyme PEPCK. A molecule of GTP is hydrolyzed to GDP during this reaction.
The next steps in the reaction are the same as reversed glycolysis. However, fructose 1,6-bisphosphatase converts fructose 1,6-bisphosphate to fructose 6-phosphate, using one water molecule and releasing one phosphate (in glycolysis, phosphofructokinase 1 converts F6P and ATP to F1,6BP and ADP). This is also the rate-limiting step of gluconeogenesis.
Glucose-6-phosphate is formed from fructose 6-phosphate by phosphoglucoisomerase (the reverse of step 2 in glycolysis). Glucose-6-phosphate can be used in other metabolic pathways or dephosphorylated to free glucose. Whereas free glucose can easily diffuse in and out of the cell, the phosphorylated form (glucose-6-phosphate) is locked in the cell, a mechanism by which intracellular glucose levels are controlled by cells.
The final step in gluconeogenesis, the formation of glucose, occurs in the lumen of the endoplasmic reticulum, where glucose-6-phosphate is hydrolyzed by glucose-6-phosphatase to produce glucose and release an inorganic phosphate. Like two steps prior, this step is not a simple reversal of glycolysis, in which hexokinase catalyzes the conversion of glucose and ATP into G6P and ADP. Glucose is shuttled into the cytoplasm by glucose transporters located in the endoplasmic reticulum's membrane.
Regulation
While most steps in gluconeogenesis are the reverse of those found in glycolysis, three regulated and strongly endergonic reactions are replaced with more kinetically favorable reactions. Hexokinase/glucokinase, phosphofructokinase, and pyruvate kinase enzymes of glycolysis are replaced with glucose-6-phosphatase, fructose-1,6-bisphosphatase, and PEP carboxykinase/pyruvate carboxylase. These enzymes are typically regulated by similar molecules, but with opposite results. For example, acetyl CoA and citrate activate gluconeogenesis enzymes (pyruvate carboxylase and fructose-1,6-bisphosphatase, respectively), while at the same time inhibiting the glycolytic enzyme pyruvate kinase. This system of reciprocal control allow glycolysis and gluconeogenesis to inhibit each other and prevents a futile cycle of synthesizing glucose to only break it down. Pyruvate kinase can be also bypassed by 86 pathways not related to gluconeogenesis, for the purpose of forming pyruvate and subsequently lactate; some of these pathways use carbon atoms originated from glucose.
The majority of the enzymes responsible for gluconeogenesis are found in the cytosol; the exceptions are mitochondrial pyruvate carboxylase and, in animals, phosphoenolpyruvate carboxykinase. The latter exists as an isozyme located in both the mitochondrion and the cytosol. The rate of gluconeogenesis is ultimately controlled by the action of a key enzyme, fructose-1,6-bisphosphatase, which is also regulated through signal transduction by cAMP and its phosphorylation.
Global control of gluconeogenesis is mediated by glucagon (released when blood glucose is low); it triggers phosphorylation of enzymes and regulatory proteins by Protein Kinase A (a cyclic AMP regulated kinase) resulting in inhibition of glycolysis and stimulation of gluconeogenesis. Insulin counteracts glucagon by inhibiting gluconeogenesis. Type 2 diabetes is marked by excess glucagon and insulin resistance from the body. Insulin can no longer inhibit the gene expression of enzymes such as PEPCK which leads to increased levels of hyperglycemia in the body. The anti-diabetic drug metformin reduces blood glucose primarily through inhibition of gluconeogenesis, overcoming the failure of insulin to inhibit gluconeogenesis due to insulin resistance.
Studies have shown that the absence of hepatic glucose production has no major effect on the control of fasting plasma glucose concentration. Compensatory induction of gluconeogenesis occurs in the kidneys and intestine, driven by glucagon, glucocorticoids, and acidosis.
Insulin resistance
In the liver, the FOX protein FOXO6 normally promotes gluconeogenesis in the fasted state, but insulin blocks FOXO6 upon feeding. In a condition of insulin resistance, insulin fails to block FOXO6 resulting in continued gluconeogenesis even upon feeding, resulting in high blood glucose (hyperglycemia).
Insulin resistance is a common feature of metabolic syndrome and type 2 diabetes. For this reason, gluconeogenesis is a target of therapy for type 2 diabetes, such as the antidiabetic drug metformin, which inhibits gluconeogenic glucose formation, and stimulates glucose uptake by cells.
Origins
Gluconeogenesis is considered one of the most ancient anabolic pathways and is likely to have been exhibited in the last universal common ancestor. Rafael F. Say and Georg Fuchs stated in 2010 that "all archaeal groups as well as the deeply branching bacterial lineages contain a bifunctional fructose 1,6-bisphosphate (FBP) aldolase/phosphatase with both FBP aldolase and FBP phosphatase activity. This enzyme is missing in most other Bacteria and in Eukaryota, and is heat-stabile even in mesophilic marine Crenarchaeota". It is proposed that fructose 1,6-bisphosphate aldolase/phosphatase was an ancestral gluconeogenic enzyme and had preceded glycolysis. But the chemical mechanisms between gluconeogenesis and glycolysis, whether it is anabolic or catabolic, are similar, suggesting they both originated at the same time. Fructose 1,6-bisphosphate is shown to be nonenzymatically synthesized continuously within a freezing solution. The synthesis is accelerated in the presence of amino acids such as glycine and lysine implying that the first anabolic enzymes were amino acids. The prebiotic reactions in gluconeogenesis can also proceed nonenzymatically at dehydration-desiccation cycles. Such chemistry could have occurred in hydrothermal environments, including temperature gradients and cycling of freezing and thawing. Mineral surfaces might have played a role in the phosphorylation of metabolic intermediates from gluconeogenesis and have to been shown to produce tetrose, hexose phosphates, and pentose from formaldehyde, glyceraldehyde, and glycolaldehyde.
| Biology and health sciences | Metabolic processes | Biology |
248808 | https://en.wikipedia.org/wiki/Algebraic%20variety | Algebraic variety | Algebraic varieties are the central objects of study in algebraic geometry, a sub-field of mathematics. Classically, an algebraic variety is defined as the set of solutions of a system of polynomial equations over the real or complex numbers. Modern definitions generalize this concept in several different ways, while attempting to preserve the geometric intuition behind the original definition.
Conventions regarding the definition of an algebraic variety differ slightly. For example, some definitions require an algebraic variety to be irreducible, which means that it is not the union of two smaller sets that are closed in the Zariski topology. Under this definition, non-irreducible algebraic varieties are called algebraic sets. Other conventions do not require irreducibility.
The fundamental theorem of algebra establishes a link between algebra and geometry by showing that a monic polynomial (an algebraic object) in one variable with complex number coefficients is determined by the set of its roots (a geometric object) in the complex plane. Generalizing this result, Hilbert's Nullstellensatz provides a fundamental correspondence between ideals of polynomial rings and algebraic sets. Using the Nullstellensatz and related results, mathematicians have established a strong correspondence between questions on algebraic sets and questions of ring theory. This correspondence is a defining feature of algebraic geometry.
Many algebraic varieties are differentiable manifolds, but an algebraic variety may have singular points while a differentiable manifold cannot. Algebraic varieties can be characterized by their dimension. Algebraic varieties of dimension one are called algebraic curves and algebraic varieties of dimension two are called algebraic surfaces.
In the context of modern scheme theory, an algebraic variety over a field is an integral (irreducible and reduced) scheme over that field whose structure morphism is separated and of finite type.
Overview and definitions
An affine variety over an algebraically closed field is conceptually the easiest type of variety to define, which will be done in this section. Next, one can define projective and quasi-projective varieties in a similar way. The most general definition of a variety is obtained by patching together smaller quasi-projective varieties. It is not obvious that one can construct genuinely new examples of varieties in this way, but Nagata gave an example of such a new variety in the 1950s.
Affine varieties
For an algebraically closed field and a natural number , let be an affine -space over , identified to through the choice of an affine coordinate system. The polynomials in the ring can be viewed as K-valued functions on by evaluating at the points in , i.e. by choosing values in K for each xi. For each set S of polynomials in , define the zero-locus Z(S) to be the set of points in on which the functions in S simultaneously vanish, that is to say
A subset V of is called an affine algebraic set if V = Z(S) for some S. A nonempty affine algebraic set V is called irreducible if it cannot be written as the union of two proper algebraic subsets. An irreducible affine algebraic set is also called an affine variety. (Some authors use the phrase affine variety to refer to any affine algebraic set, irreducible or not.)
Affine varieties can be given a natural topology by declaring the closed sets to be precisely the affine algebraic sets. This topology is called the Zariski topology.
Given a subset V of , we define I(V) to be the ideal of all polynomial functions vanishing on V:
For any affine algebraic set V, the coordinate ring or structure ring of V is the quotient of the polynomial ring by this ideal.
Projective varieties and quasi-projective varieties
Let be an algebraically closed field and let be the projective n-space over . Let in be a homogeneous polynomial of degree d. It is not well-defined to evaluate on points in in homogeneous coordinates. However, because is homogeneous, meaning that , it does make sense to ask whether vanishes at a point . For each set S of homogeneous polynomials, define the zero-locus of S to be the set of points in on which the functions in S vanish:
A subset V of is called a projective algebraic set if V = Z(S) for some S. An irreducible projective algebraic set is called a projective variety.
Projective varieties are also equipped with the Zariski topology by declaring all algebraic sets to be closed.
Given a subset V of , let I(V) be the ideal generated by all homogeneous polynomials vanishing on V. For any projective algebraic set V, the coordinate ring of V is the quotient of the polynomial ring by this ideal.
A quasi-projective variety is a Zariski open subset of a projective variety. Notice that every affine variety is quasi-projective. Notice also that the complement of an algebraic set in an affine variety is a quasi-projective variety; in the context of affine varieties, such a quasi-projective variety is usually not called a variety but a constructible set.
Abstract varieties
In classical algebraic geometry, all varieties were by definition quasi-projective varieties, meaning that they were open subvarieties of closed subvarieties of a projective space. For example, in Chapter 1 of Hartshorne a variety over an algebraically closed field is defined to be a quasi-projective variety, but from Chapter 2 onwards, the term variety (also called an abstract variety) refers to a more general object, which locally is a quasi-projective variety, but when viewed as a whole is not necessarily quasi-projective; i.e. it might not have an embedding into projective space. So classically the definition of an algebraic variety required an embedding into projective space, and this embedding was used to define the topology on the variety and the regular functions on the variety. The disadvantage of such a definition is that not all varieties come with natural embeddings into projective space. For example, under this definition, the product is not a variety until it is embedded into a larger projective space; this is usually done by the Segre embedding. Furthermore, any variety that admits one embedding into projective space admits many others, for example by composing the embedding with the Veronese embedding; thus many notions that should be intrinsic, such as that of a regular function, are not obviously so.
The earliest successful attempt to define an algebraic variety abstractly, without an embedding, was made by André Weil. In his Foundations of Algebraic Geometry, using valuations. Claude Chevalley made a definition of a scheme, which served a similar purpose, but was more general. However, Alexander Grothendieck's definition of a scheme is more general still and has received the most widespread acceptance. In Grothendieck's language, an abstract algebraic variety is usually defined to be an integral, separated scheme of finite type over an algebraically closed field, although some authors drop the irreducibility or the reducedness or the separateness condition or allow the underlying field to be not algebraically closed. Classical algebraic varieties are the quasiprojective integral separated finite type schemes over an algebraically closed field.
Existence of non-quasiprojective abstract algebraic varieties
One of the earliest examples of a non-quasiprojective algebraic variety were given by Nagata. Nagata's example was not complete (the analog of compactness), but soon afterwards he found an algebraic surface that was complete and non-projective. Since then other examples have been found: for example, it is straightforward to construct toric varieties that are not quasi-projective but complete.
Examples
Subvariety
A subvariety is a subset of a variety that is itself a variety (with respect to the topological structure induced by the ambient variety). For example, every open subset of a variety is a variety. | Mathematics | Algebra | null |
248860 | https://en.wikipedia.org/wiki/Theodolite | Theodolite | A theodolite () is a precision optical instrument for measuring angles between designated visible points in the horizontal and vertical planes. The traditional use has been for land surveying, but it is also used extensively for building and infrastructure construction, and some specialized applications such as meteorology and rocket launching.
It consists of a moveable telescope mounted so it can rotate around horizontal and vertical axes and provide angular readouts. These indicate the orientation of the telescope, and are used to relate the first point sighted through the telescope to subsequent sightings of other points from the same theodolite position. These angles can be measured with accuracies down to microradians or seconds of arc. From these readings a plan can be drawn, or objects can be positioned in accordance with an existing plan. The modern theodolite has evolved into what is known as a total station where angles and distances are measured electronically, and are read directly to computer memory.
In a transit theodolite, the telescope is short enough to rotate about the trunnion axis, turning the telescope through the vertical plane through the zenith; for non-transit instruments vertical rotation is restricted to a limited arc.
The optical level is sometimes mistaken for a theodolite, but it does not measure vertical angles, and is used only for leveling on a horizontal plane (though often combined with medium accuracy horizontal range and direction measurements).
Principles of operation
Preparation for making sightings
Temporary adjustments are a set of operations necessary in order to make a theodolite ready for taking observations at a station. These include its setting up, centering, leveling up and elimination of parallax, and are achieved in four steps:
Setting up: fixing the theodolite onto a tripod along with approximate leveling and centering over the station mark.
Centering: bringing the vertical axis of theodolite immediately over station mark using a centering plate also known as a tribrach.
Leveling: leveling of the base of the instrument to make the vertical axis vertical usually with an in-built bubble-level.
Focusing: removing parallax error by proper focusing of objective and eye-piece. The eye-piece only requires adjustment once at a station. The objective will be re-focused for each subsequent sighting from this station because of the different distances to the target.
Sightings
Sightings are taken by the surveyor, who adjusts the telescope's vertical and horizontal angular orientation so the cross-hairs align with the desired sighting point. Both angles are read either from exposed or internal scales and recorded. The next object is then sighted and recorded without moving the position of the instrument and tripod.
The earliest angular readouts were from open vernier scales directly visible to the eye. Gradually these scales were enclosed for physical protection, and finally became an indirect optical readout, with convoluted light paths to bring them to a convenient place on the instrument for viewing. The modern digital theodolites have electronic displays.
Errors in measurement
Index error The angles in the vertical axis should read 90° (100 grad) when the sight axis is horizontal, or 270° (300 grad) when the instrument is transited. Half of the difference between the two positions is called the index error. This can only be checked on transit instruments.
Horizontal axis error The horizontal and vertical axes of a theodolite must be perpendicular; if not then a horizontal axis error exists. This can be tested by aligning the tubular spirit bubble parallel to a line between two footscrews and setting the bubble central. A horizontal axis error is present if the bubble runs off central when the tubular spirit bubble is reversed (turned through 180°). To adjust, the operator removes half the amount the bubble has run off using the adjusting screw, then re-level, test and refine the adjustment.
Collimation error The optical axis of the telescope must also be perpendicular to the horizontal axis; if not, then a collimation error exists.
Index error, horizontal-axis error (trunnion-axis error) and collimation error are regularly determined by calibration and are removed by mechanical adjustment. Their existence is taken into account in the choice of measurement procedure in order to eliminate their effect on the measurement results of the theodolite.
History
Historical background
Prior to the theodolite, instruments such as the groma, geometric square and the dioptra, and various other graduated circles (see circumferentor) and semicircles (see graphometer) were used to obtain either vertical or horizontal angle measurements. Over time their functions were combined into a single instrument that could measure both angles simultaneously.
The first occurrence of the word "theodolite" is found in the surveying textbook A geometric practice named Pantometria (1571) by Leonard Digges. The origin of the word is unknown. The first part of the Neo-Latin theo-delitus might stem from the Greek , "to behold or look attentively upon" The second part is often attributed to an unscholarly variation of the Greek word: , meaning "evident" or "clear". Other Neo-Latin or Greek derivations have been suggested as well as an English origin from "the alidade".
The early forerunners of the theodolite were sometimes azimuth instruments for measuring horizontal angles, while others had an altazimuth mount for measuring horizontal and vertical angles. Gregorius Reisch illustrated an altazimuth instrument in the appendix of his 1512 book Margarita Philosophica. Martin Waldseemüller, a topographer and cartographer made the device in that year calling it the polimetrum. In Digges's book of 1571, the term "theodolite" was applied to an instrument for measuring horizontal angles only, but he also described an instrument that measured both altitude and azimuth which he called a instrument . Possibly the first instrument approximating to a true theodolite was the built by Josua Habemel in 1576, complete with compass and tripod. The 1728 Cyclopaedia compares "graphometer" to "half-theodolite". As late as the 19th century, the instrument for measuring horizontal angles only was called a simple theodolite and the altazimuth instrument, the plain theodolite.
The first instrument to combine the essential features of the modern theodolite was built in 1725 by Jonathan Sisson. This instrument had an altazimuth mount with a sighting telescope. The base plate had spirit levels, compass and adjusting screws. The circles were read with a vernier scale.
Development of the theodolite
The theodolite became a modern, accurate instrument in 1787, with the introduction of Jesse Ramsden's famous great theodolite, which he created using a very accurate dividing engine of his own design. Ramsden's instruments were used for the Principal Triangulation of Great Britain. At this time the highest precision instruments were made in England by such makers as Edward Troughton. Later the first practical German theodolites were made by Breithaupt together with Utzschneider, Reichenbach and Fraunhofer.
As technology progressed the vertical partial circle was replaced with a full circle, and both vertical and horizontal circles were finely graduated. This was the transit theodolite. This type of theodolite was developed from 18th century astronomical Transit instruments used to measure accurate star positions. The technology was transferred to theodolites in the early 19th century by instrument makers such as Edward Troughton and William Simms and became the standard theodolite design. Development of the theodolite was spurred on by specific needs. In the 1820s progress on national surveying projects such as the Ordnance Survey in Britain produced a requirement for theodolites capable of providing sufficient accuracy for large scale triangulation and mapping. The Survey of India at this time produced a requirement for more rugged and stable instruments such as the Everest pattern theodolite with its lower center of gravity.
Railway engineers working in the 1830s in Britain commonly referred to a theodolite as a "Transit". The 1840s was the start of a period of rapid railway building in many parts of the world which resulted in a high demand for theodolites wherever railways were being constructed. It was also popular with American railroad engineers pushing west, and it replaced the railroad compass, sextant and octant. Theodolites were later adapted to a wider variety of mountings and uses. In the 1870s, an interesting waterborne version of the theodolite (using a pendulum device to counteract wave movement) was invented by Edward Samuel Ritchie. It was used by the U.S. Navy to take the first precision surveys of American harbors on the Atlantic and Gulf coasts.
In the early 1920s a step change in theodolite design occurred with the introduction of the Wild T2 made by the Swiss Wild Heerbrugg company. Heinrich Wild designed a theodolite with divided glass circles with readings from both sides presented at a single eyepiece close to the telescope so the observer did not have to move to read them. The Wild instruments were not only smaller, easier to use and more accurate than contemporary rivals but also sealed from rain and dust. Canadian surveyors reported that while the Wild T2 with 3.75 inch circles was not able to provide the accuracy for primary triangulation it was the equal in accuracy to a 12 inch traditional design. The Wild T2, T3, and A1 instruments were made for many years.
In 1926 a conference was held at Tavistock in Devon, UK where Wild theodolites were compared with British ones. The Wild product outclassed the British theodolites so manufacturers such as Cooke, Troughton & Simms and Hilger & Watts set about improving the accuracy of their products to match their competition. Cooke, Troughton and Simms developed the Tavistock pattern theodolite and later the Vickers V. 22.
Wild went on to develop the DK1, DKM1, DM2, DKM2, and DKM3 for Kern Aarau company. With continuing refinements, instruments steadily evolved into the modern theodolite used by surveyors today. By 1977 Wild, Kern and Hewlett-Packard were all offering "Total stations" which combined angular measurements, electronic distance measurement and microchip functions in a single unit.
Operation in surveying
Triangulation, as invented by Gemma Frisius around 1533, consists of making such direction plots of the surrounding landscape from two separate standpoints. The two graphing papers are superimposed, providing a scale model of the landscape, or rather the targets in it. The true scale can be obtained by measuring one distance both in the real terrain and in the graphical representation.
Modern triangulation as, e.g., practiced by Snellius, is the same procedure executed by numerical means. Photogrammetric block adjustment of stereo pairs of aerial photographs is a modern, three-dimensional variant.
In the late 1780s, Jesse Ramsden, a Yorkshireman from Halifax, England who had developed the dividing engine for dividing angular scales accurately to within a second of arc (≈ 0.0048 mrad or 4.8 μrad), was commissioned to build a new instrument for the British Ordnance Survey. The Ramsden theodolite was used over the next few years to map the whole of southern Britain by triangulation.
In network measurement, the use of forced centering speeds up operations while maintaining the highest precision. The theodolite or the target can be rapidly removed from, or socketed into, the forced centering plate with sub-millimeter precision. Nowadays GPS antennas used for geodetic positioning use a similar mounting system. The height of the reference point of the theodolite—or the target—above the ground benchmark must be measured precisely.
Transit theodolite
The term transit theodolite, or transit for short, refers to a type of theodolite where the telescope is short enough to rotate in a full circle on its horizontal axis as well as around its vertical axis. It features a vertical circle which is graduated through the full 360 degrees and a telescope that could "flip over" ("transit the scope"). By reversing the telescope and at the same time rotating the instrument through 180 degrees about the vertical axis, the instrument can be used in 'plate-left' or 'plate-right' modes ('plate' refers to the vertical protractor circle). By measuring the same horizontal and vertical angles in these two modes and then averaging the results, centering and collimating errors in the instrument can be eliminated. Some transit instruments are capable of reading angles directly to thirty arc-seconds (≈ 0.15 mrad). Modern theodolites are usually of the transit-theodolite design, but engraved plates have been replaced with glass plates designed to be read with light-emitting diodes and computer circuitry, greatly improving accuracy up to arc-second (≈ 0.005 mrad) levels.
Use with weather balloons
There is a long history of theodolite use in measuring winds aloft, by using specially-manufactured theodolites to track the horizontal and vertical angles of special weather balloons called ceiling balloons or pilot balloons (pibal). Early attempts at this were made in the opening years of the nineteenth century, but the instruments and procedures weren't fully developed until a hundred years later. This method was extensively used in World War II and thereafter, and was gradually replaced by radio and GPS measuring systems from the 1980s onward.
The pibal theodolite uses a prism to bend the optical path by 90 degrees so the operator's eye position does not change as the elevation is changed through a complete 180 degrees. The theodolite is typically mounted on a rugged steel stand, set up so it is level and pointed north, with the altitude and azimuth scales reading zero degrees. A balloon is released in front of the theodolite, and its position is precisely tracked, usually once a minute. The balloons are carefully constructed and filled, so their rate of ascent can be known fairly accurately in advance. Mathematical calculations on time, rate of ascent, azimuth and angular altitude can produce good estimates of wind speed and direction at various altitudes.
Modern electronic theodolites
In modern electronic theodolites, the readout of the horizontal and vertical circles is usually done with a rotary encoder. These produce signals indicating the altitude and azimuth of the telescope which are fed to a microprocessor. CCD sensors have been added to the focal plane of the telescope allowing both auto-targeting and the automated measurement of residual target offset. All this is implemented in embedded software of the processor.
Many modern theodolites are equipped with integrated electro-optical distance measuring devices, generally infrared based, allowing the measurement in one step of complete three-dimensional vectors—albeit in instrument-defined polar coordinates, which can then be transformed to a preexisting coordinate system in the area by means of a sufficient number of control points. This technique is called a resection solution or free station position surveying and is widely used in mapping surveying.
Such instruments are "intelligent" theodolites called self-registering tacheometers or colloquially "total stations", and perform all the necessary angular and distance calculations, and the results or raw data can be downloaded to external processors, such as ruggedized laptops, PDAs or programmable calculators.
Gyrotheodolites
A gyrotheodolite is used when the north-south reference bearing of the meridian is required in the absence of astronomical star sights. This occurs mainly in the underground mining industry and in tunnel engineering. For example, where a conduit must pass under a river, a vertical shaft on each side of the river might be connected by a horizontal tunnel. A gyrotheodolite can be operated at the surface and then again at the foot of the shafts to identify the directions needed to tunnel between the base of the two shafts. Unlike an artificial horizon or inertial navigation system, a gyrotheodolite cannot be relocated while it is operating. It must be restarted again at each site.
The gyrotheodolite comprises a normal theodolite with an attachment that contains a gyrocompass, a device which senses the rotation of the Earth in order to find true north and thus, in conjunction with the direction of gravity, the plane of the meridian. The meridian is the plane that contains both the axis of the Earth's rotation and the observer. The intersection of the meridian plane with the horizontal defines the true north-south direction found in this way. Unlike magnetic compasses, gyrocompasses are able to find true north, the surface direction toward the north pole.
A gyrotheodolite will function at the equator and in both the northern and southern hemispheres. The meridian is undefined at the geographic poles. A gyrotheodolite cannot be used at the poles where the Earth's axis is precisely perpendicular to the horizontal axis of the spinner, indeed it is not normally used within about 15 degrees of the pole where the angle between the earth's rotation and the direction of gravity is too small for it to work reliably. When available, astronomical star sights are able to give the meridian bearing to better than one hundred times the accuracy of the gyrotheodolite. Where this extra precision is not required, the gyrotheodolite is able to produce a result quickly without the need for night observations.
| Technology | Optical instruments | null |
248932 | https://en.wikipedia.org/wiki/Software%20development | Software development | Software development is the process of designing and implementing a software solution to satisfy a user. The process is more encompassing than programming, writing code, in that it includes conceiving the goal, evaluating feasibility, analyzing requirements, design, testing and release. The process is part of software engineering which also includes organizational management, project management, configuration management and other aspects.
Software development involves many skills and job specializations including programming, testing, documentation, graphic design, user support, marketing, and fundraising.
Software development involves many tools including: compiler, integrated development environment (IDE), version control, computer-aided software engineering, and word processor.
The details of the process used for a development effort varies. The process may be confined to a formal, documented standard, or it can be customized and emergent for the development effort. The process may be sequential, in which each major phase (i.e. design, implement and test) is completed before the next begins, but an iterative approach where small aspects are separately designed, implemented and tested can reduce risk and cost and increase quality.
Methodologies
Each of the available methodologies are best suited to specific kinds of projects, based on various technical, organizational, project, and team considerations.
The simplest methodology is the "code and fix", typically used by a single programmer working on a small project. After briefly considering the purpose of the program, the programmer codes it and runs it to see if it works. When they are done, the product is released. This methodology is useful for prototypes but cannot be used for more elaborate programs.
In the top-down waterfall model, feasibility, analysis, design, development, quality assurance, and implementation occur sequentially in that order. This model requires one step to be complete before the next begins, causing delays, and makes it impossible to revise previous steps if necessary.
With iterative processes these steps are interleaved with each other for improved flexibility, efficiency, and more realistic scheduling. Instead of completing the project all at once, one might go through most of the steps with one component at a time. Iterative development also lets developers prioritize the most important features, enabling lower priority ones to be dropped later on if necessary. Agile is one popular method, originally intended for small or medium sized projects, that focuses on giving developers more control over the features that they work on to reduce the risk of time or cost overruns. Derivatives of agile include extreme programming and Scrum. Open-source software development typically uses agile methodology with concurrent design, coding, and testing, due to reliance on a distributed network of volunteer contributors.
Beyond agile, some companies integrate information technology (IT) operations with software development, which is called DevOps or DevSecOps including computer security. DevOps includes continuous development, testing, integration of new code in the version control system, deployment of the new code, and sometimes delivery of the code to clients. The purpose of this integration is to deliver IT services more quickly and efficiently.
Another focus in many programming methodologies is the idea of trying to catch issues such as security vulnerabilities and bugs as early as possible (shift-left testing) to reduce the cost of tracking and fixing them.
In 2009, it was estimated that 32 percent of software projects were delivered on time and budget, and with the full functionality. An additional 44 percent were delivered, but missing at least one of these features. The remaining 24 percent were cancelled prior to release.
Steps
Software development life cycle refers to the systematic process of developing applications.
Feasibility
The sources of ideas for software products are plentiful. These ideas can come from market research including the demographics of potential new customers, existing customers, sales prospects who rejected the product, other internal software development staff, or a creative third party. Ideas for software products are usually first evaluated by marketing personnel for economic feasibility, fit with existing channels of distribution, possible effects on existing product lines, required features, and fit with the company's marketing objectives. In the marketing evaluation phase, the cost and time assumptions become evaluated. The feasibility analysis estimates the project's return on investment, its development cost and timeframe. Based on this analysis, the company can make a business decision to invest in further development. After deciding to develop the software, the company is focused on delivering the product at or below the estimated cost and time, and with a high standard of quality (i.e., lack of bugs) and the desired functionality. Nevertheless, most software projects run late and sometimes compromises are made in features or quality to meet a deadline.
Analysis
Software analysis begins with a requirements analysis to capture the business needs of the software. Challenges for the identification of needs are that current or potential users may have different and incompatible needs, may not understand their own needs, and change their needs during the process of software development. Ultimately, the result of analysis is a detailed specification for the product that developers can work from. Software analysts often decompose the project into smaller objects, components that can be reused for increased cost-effectiveness, efficiency, and reliability. Decomposing the project may enable a multi-threaded implementation that runs significantly faster on multiprocessor computers.
During the analysis and design phases of software development, structured analysis is often used to break down the customer's requirements into pieces that can be implemented by software programmers. The underlying logic of the program may be represented in data-flow diagrams, data dictionaries, pseudocode, state transition diagrams, and/or entity relationship diagrams. If the project incorporates a piece of legacy software that has not been modeled, this software may be modeled to help ensure it is correctly incorporated with the newer software.
Design
Design involves choices about the implementation of the software, such as which programming languages and database software to use, or how the hardware and network communications will be organized. Design may be iterative with users consulted about their needs in a process of trial and error. Design often involves people expert in aspect such as database design, screen architecture, and the performance of servers and other hardware. Designers often attempt to find patterns in the software's functionality to spin off distinct modules that can be reused with object-oriented programming. An example of this is the model–view–controller, an interface between a graphical user interface and the backend.
Programming
The central feature of software development is creating and understanding the software that implements the desired functionality. There are various strategies for writing the code. Cohesive software has various components that are independent from each other. Coupling is the interrelation of different software components, which is viewed as undesirable because it increases the difficulty of maintenance. Often, software programmers do not follow industry best practices, resulting in code that is inefficient, difficult to understand, or lacking documentation on its functionality. These standards are especially likely to break down in the presence of deadlines. As a result, testing, debugging, and revising the code becomes much more difficult. Code refactoring, for example adding more comments to the code, is a solution to improve the understandability of code.
Testing
Testing is the process of ensuring that the code executes correctly and without errors. Debugging is performed by each software developer on their own code to confirm that the code does what it is intended to. In particular, it is crucial that the software executes on all inputs, even if the result is incorrect. Code reviews by other developers are often used to scrutinize new code added to the project, and according to some estimates dramatically reduce the number of bugs persisting after testing is complete. Once the code has been submitted, quality assurance—a separate department of non-programmers for most large companies—test the accuracy of the entire software product. Acceptance tests derived from the original software requirements are a popular tool for this. Quality testing also often includes stress and load checking (whether the software is robust to heavy levels of input or usage), integration testing (to ensure that the software is adequately integrated with other software), and compatibility testing (measuring the software's performance across different operating systems or browsers). When tests are written before the code, this is called test-driven development.
Production
Production is the phase in which software is deployed to the end user. During production, the developer may create technical support resources for users or a process for fixing bugs and errors that were not caught earlier. There might also be a return to earlier development phases if user needs changed or were misunderstood.
Workers
Software development is performed by software developers, usually working on a team. Efficient communications between team members is essential to success. This is more easily achieved if the team is small, used to working together, and located near each other. Communications also help identify problems at an earlier state of development and avoid duplicated effort. Many development projects avoid the risk of losing essential knowledge held by only one employee by ensuring that multiple workers are familiar with each component. Software development involves professionals from various fields, not just software programmers but also individuals specialized in testing, documentation writing, graphic design, user support, marketing, and fundraising. Although workers for proprietary software are paid, most contributors to open-source software are volunteers. Alternately, they may be paid by companies whose business model does not involve selling the software, but something else—such as services and modifications to open source software.
Models and tools
Computer-aided software engineering
Computer-aided software engineering (CASE) is tools for the partial automation of software development. CASE enables designers to sketch out the logic of a program, whether one to be written, or an already existing one to help integrate it with new code or reverse engineer it (for example, to change the programming language).
Documentation
Documentation comes in two forms that are usually kept separate—that intended for software developers, and that made available to the end user to help them use the software. Most developer documentation is in the form of code comments for each file, class, and method that cover the application programming interface (API)—how the piece of software can be accessed by another—and often implementation details. This documentation is helpful for new developers to understand the project when they begin working on it. In agile development, the documentation is often written at the same time as the code. User documentation is more frequently written by technical writers.
Effort estimation
Accurate estimation is crucial at the feasibility stage and in delivering the product on time and within budget. The process of generating estimations is often delegated by the project manager. Because the effort estimation is directly related to the size of the complete application, it is strongly influenced by addition of features in the requirements—the more requirements, the higher the development cost. Aspects not related to functionality, such as the experience of the software developers and code reusability, are also essential to consider in estimation. , most of the tools for estimating the amount of time and resources for software development were designed for conventional applications and are not applicable to web applications or mobile applications.
Integrated development environment
An integrated development environment (IDE) supports software development with enhanced features compared to a simple text editor. IDEs often include automated compiling, syntax highlighting of errors, debugging assistance, integration with version control, and semi-automation of tests.
Version control
Version control is a popular way of managing changes made to the software. Whenever a new version is checked in, the software saves a backup of all modified files. If multiple programmers are working on the software simultaneously, it manages the merging of their code changes. The software highlights cases where there is a conflict between two sets of changes and allows programmers to fix the conflict.
View model
A view model is a framework that provides the viewpoints on the system and its environment, to be used in the software development process. It is a graphical representation of the underlying semantics of a view.
The purpose of viewpoints and views is to enable human engineers to comprehend very complex systems and to organize the elements of the problem around domains of expertise. In the engineering of physically intensive systems, viewpoints often correspond to capabilities and responsibilities within the engineering organization.
Fitness functions
Fitness functions are automated and objective tests to ensure that the new developments don't deviate from the established constraints, checks and compliance controls.
Intellectual property
Intellectual property can be an issue when developers integrate open-source code or libraries into a proprietary product, because most open-source licenses used for software require that modifications be released under the same license. As an alternative, developers may choose a proprietary alternative or write their own software module.
| Technology | Software development: General | null |
248988 | https://en.wikipedia.org/wiki/Product%20rule | Product rule | In calculus, the product rule (or Leibniz rule or Leibniz product rule) is a formula used to find the derivatives of products of two or more functions. For two functions, it may be stated in Lagrange's notation as or in Leibniz's notation as
The rule may be extended or generalized to products of three or more functions, to a rule for higher-order derivatives of a product, and to other contexts.
Discovery
Discovery of this rule is credited to Gottfried Leibniz, who demonstrated it using "infinitesimals" (a precursor to the modern differential). (However, J. M. Child, a translator of Leibniz's papers, argues that it is due to Isaac Barrow.) Here is Leibniz's argument: Let u and v be functions. Then d(uv) is the same thing as the difference between two successive uv'''s; let one of these be uv, and the other u+du times v+dv; then:
Since the term du·dv is "negligible" (compared to du and dv), Leibniz concluded that
and this is indeed the differential form of the product rule. If we divide through by the differential dx, we obtain
which can also be written in Lagrange's notation as
Examples
Suppose we want to differentiate By using the product rule, one gets the derivative (since the derivative of is and the derivative of the sine function is the cosine function).
One special case of the product rule is the constant multiple rule, which states: if is a number, and is a differentiable function, then is also differentiable, and its derivative is This follows from the product rule since the derivative of any constant is zero. This, combined with the sum rule for derivatives, shows that differentiation is linear.
The rule for integration by parts is derived from the product rule, as is (a weak version of) the quotient rule. (It is a "weak" version in that it does not prove that the quotient is differentiable but only says what its derivative is it is differentiable.)
Proofs
Limit definition of derivative
Let and suppose that and are each differentiable at . We want to prove that is differentiable at and that its derivative, , is given by . To do this, (which is zero, and thus does not change the value) is added to the numerator to permit its factoring, and then properties of limits are used.
The fact that follows from the fact that differentiable functions are continuous.
Linear approximations
By definition, if are differentiable at , then we can write linear approximations:
and
where the error terms are small with respect to h: that is, also written . Then:
The "error terms" consist of items such as and which are easily seen to have magnitude Dividing by and taking the limit gives the result.
Quarter squares
This proof uses the chain rule and the quarter square function with derivative . We have:
and differentiating both sides gives:
Multivariable chain rule
The product rule can be considered a special case of the chain rule for several variables, applied to the multiplication function :
Non-standard analysis
Let u and v be continuous functions in x, and let dx, du and dv be infinitesimals within the framework of non-standard analysis, specifically the hyperreal numbers. Using st to denote the standard part function that associates to a finite hyperreal number the real infinitely close to it, this gives
This was essentially Leibniz's proof exploiting the transcendental law of homogeneity (in place of the standard part above).
Smooth infinitesimal analysis
In the context of Lawvere's approach to infinitesimals, let be a nilsquare infinitesimal. Then and , so that
since Dividing by then gives or .
Logarithmic differentiation
Let . Taking the absolute value of each function and the natural log of both sides of the equation,
Applying properties of the absolute value and logarithms,
Taking the logarithmic derivative of both sides and then solving for :
Solving for and substituting back for gives:
Note: Taking the absolute value of the functions is necessary for the logarithmic differentiation of functions that may have negative values, as logarithms are only real-valued for positive arguments. This works because , which justifies taking the absolute value of the functions for logarithmic differentiation.
Generalizations
Product of more than two factors
The product rule can be generalized to products of more than two factors. For example, for three factors we have
For a collection of functions , we have
The logarithmic derivative provides a simpler expression of the last form, as well as a direct proof that does not involve any recursion. The logarithmic derivative of a function , denoted here , is the derivative of the logarithm of the function. It follows that
Using that the logarithm of a product is the sum of the logarithms of the factors, the sum rule for derivatives gives immediately
The last above expression of the derivative of a product is obtained by multiplying both members of this equation by the product of the
Higher derivatives
It can also be generalized to the general Leibniz rule for the nth derivative of a product of two factors, by symbolically expanding according to the binomial theorem:
Applied at a specific point x, the above formula gives:
Furthermore, for the nth derivative of an arbitrary number of factors, one has a similar formula with multinomial coefficients:
Higher partial derivatives
For partial derivatives, we have
where the index runs through all subsets of , and is the cardinality of . For example, when ,
Banach space
Suppose X, Y, and Z are Banach spaces (which includes Euclidean space) and B : X × Y → Z is a continuous bilinear operator. Then B is differentiable, and its derivative at the point (x,y) in X × Y is the linear map D(x,y)B : X × Y → Z given by
This result can be extended to more general topological vector spaces.
In vector calculus
The product rule extends to various product operations of vector functions on :
For scalar multiplication:
For dot product:
For cross product of vector functions on :
There are also analogues for other analogs of the derivative: if f and g are scalar fields then there is a product rule with the gradient:
Such a rule will hold for any continuous bilinear product operation. Let B : X × Y → Z be a continuous bilinear map between vector spaces, and let f and g be differentiable functions into X and Y, respectively. The only properties of multiplication used in the proof using the limit definition of derivative is that multiplication is continuous and bilinear. So for any continuous bilinear operation,
This is also a special case of the product rule for bilinear maps in Banach space.
Derivations in abstract algebra and differential geometry
In abstract algebra, the product rule is the defining property of a derivation. In this terminology, the product rule states that the derivative operator is a derivation on functions.
In differential geometry, a tangent vector to a manifold M at a point p may be defined abstractly as an operator on real-valued functions which behaves like a directional derivative at p: that is, a linear functional v which is a derivation,
Generalizing (and dualizing) the formulas of vector calculus to an n-dimensional manifold M, one may take differential forms of degrees k and l, denoted , with the wedge or exterior product operation , as well as the exterior derivative . Then one has the graded Leibniz rule:
Applications
Among the applications of the product rule is a proof that
when n is a positive integer (this rule is true even if n is not positive or is not an integer, but the proof of that must rely on other methods). The proof is by mathematical induction on the exponent n. If n = 0 then xn is constant and nxn − 1 = 0. The rule holds in that case because the derivative of a constant function is 0. If the rule holds for any particular exponent n, then for the next value, n + 1, we have
Therefore, if the proposition is true for n, it is true also for n + 1, and therefore for all natural n''.
| Mathematics | Differential calculus | null |
249033 | https://en.wikipedia.org/wiki/Multiplication%20sign | Multiplication sign | The multiplication sign (), also known as the times sign or the dimension sign, is a mathematical symbol used to denote the operation of multiplication, which results in a product.
The symbol is also used in botany, in botanical hybrid names.
The form is properly a four-fold rotationally symmetric saltire. The multiplication sign is similar to a lowercase X () which is not a four-fold rotationally symmetric saltire.
History
The earliest known use of the symbol to indicate multiplication appears in an anonymous appendix to the 1618 edition of John Napier's . This appendix has been attributed to William Oughtred, who used the same symbol in his 1631 algebra text, , stating:Multiplication of species [i.e. unknowns] connects both proposed magnitudes with the symbol 'in' or : or ordinarily without the symbol if the magnitudes be denoted with one letter.
Other works have been identified in which crossed diagonals appear in diagrams involving multiplied numbers, such as Robert Recorde's The Ground of Arts and Oswald Schreckenfuchs's 1551 edition of Almagest, but these are not symbolizations.
Uses
In mathematics, the symbol × has a number of uses, including
Multiplication of two numbers, where it is read as "times" or "multiplied by"
Cross product of two vectors, where it is usually read as "cross"
Cartesian product of two sets, where it is usually read as "cross"
Geometric dimension of an object, such as noting that a room is 10 feet × 12 feet in area, where it is usually read as "by" (e.g., "10 feet by 12 feet")
Screen resolution in pixels, such as 1920 pixels across × 1080 pixels down. Read as "by".
Dimensions of a matrix, where it is usually read as "by"
A statistical interaction between two explanatory variables, where it is usually read as "by"
In biology, the multiplication sign is used in a botanical hybrid name, for instance Ceanothus papillosus × impressus (a hybrid between C. papillosus and C. impressus) or Crocosmia × crocosmiiflora (a hybrid between two other species of Crocosmia). However, the communication of these hybrid names with a Latin letter "x" is common, when the actual "×" symbol is not readily available.
The multiplication sign is also used by historians for an event between two dates. When employed between two dates for example 1225 and 1232 the expression "1225×1232" means "no earlier than 1225 and no later than 1232".
A monadic symbol is used by the APL programming language to denote the sign function.
Similar notations
The lower-case Latin letter is sometimes used in place of the multiplication sign. This is considered incorrect in mathematical writing.
In algebraic notation, widely used in mathematics, a multiplication symbol is usually omitted wherever it would not cause confusion: " multiplied by " can be written as or .
Other symbols can also be used to denote multiplication, often to reduce confusion between the multiplication sign × and the common variable . In some countries, such as Germany, the primary symbol for multiplication is the "dot operator" (as in ). This symbol is also used in compound units of measurement, e.g., N⋅m (see International System of Units#Lexicographic conventions). In algebra, it is a notation to resolve ambiguity (for instance, "b times 2" may be written as , to avoid being confused with a value called ). This notation is used wherever multiplication should be written explicitly, such as in " for "; this usage is also seen in English-language texts. In some languages, the use of full stop as a multiplication symbol, such as , is common when the symbol for decimal point is comma.
Historically, computer language syntax was restricted to the ASCII character set, and the asterisk became the de facto symbol for the multiplication operator. This selection is reflected in the numeric keypad on English-language keyboards, where the arithmetic operations of addition, subtraction, multiplication and division are represented by the keys , , and , respectively.
Typing the character
Unicode and HTML entities
Other variants and related characters:
(a zero-width space indicating multiplication)
(the interpunct, may be easier to type than the dot operator)
(intended to explicitly denote the cross product of two vectors)
| Mathematics | Basics | null |
249116 | https://en.wikipedia.org/wiki/Liquorice | Liquorice | Liquorice (Commonwealth English) or licorice (American English; see spelling differences; ) is the common name of Glycyrrhiza glabra, a flowering plant of the bean family Fabaceae, from the root of which a sweet, aromatic flavouring is extracted.
The liquorice plant is an herbaceous perennial legume native to West Asia, North Africa, and Southern Europe. Liquorice is used as a flavouring in confectionery, tobacco, beverages, and pharmaceuticals, and is marketed as a dietary supplement.
Liquorice extracts have been used in herbalism and traditional medicine. Excessive consumption of liquorice (more than per day of pure glycyrrhizinic acid, a key component of liquorice) can lead to undesirable consequences. Clinically, it is suspected that overindulgence in liquorice may manifest as unexplained hypertension, low blood potassium levels (hypokalemia), and muscle weakness in individuals. Consuming liquorice should be avoided during pregnancy.
Etymology
The word liquorice (UK, CAN), or licorice (US), is derived via the Anglo-French , from Late Latin , itself ultimately derived from Greek (the Modern Greek spelling of the genus is ) literally meaning 'sweet root' and referring to Glycyrrhiza glabra.
The latter gives the plant binomial name with glabra meaning smooth and referring to the plant's smooth husks; the former came to being via the influence of , 'to become fluid', reflecting the method of extracting the sweet component from the roots.
, its English common name is spelled 'liquorice' in most of the Commonwealth, but 'licorice' is also used in some countries.
Description
Liquorice is a herbaceous perennial, growing to in height, with pinnate leaves about long, with 9–17 leaflets. The flowers are long, purple to pale whitish blue, produced in a loose inflorescence. The fruit is an oblong pod, long, containing several seeds. The roots are stoloniferous.
Chemistry
Liquorice root contains triterpenoids, polyphenols, and polysaccharides. Flavonoids account for the yellow root color. The principal glycoside, glycyrrhizin, exists in content of 7% to 10%, depending on cultivation practices. The isoflavene glabrene and the isoflavane glabridin, found in the roots of liquorice, are phytoestrogens.
The scent of liquorice root comes from a complex and variable combination of compounds, of which anethole is some 3% of total volatiles. Much of the sweetness in liquorice comes from glycyrrhizin, which has 30–50 times the sweetness of sugar. The sweetness is different from sugar, being less instant, tart, and lasting longer.
Cultivation and uses
Liquorice grows best in well-drained soils in deep valleys with full sun. It is harvested in the autumn two to three years after planting. Countries producing liquorice include Turkey, Greece, Iran, and Iraq.
Tobacco
Liquorice is used as a flavouring agent for tobacco, for flavour-enhancing and moistening agents in the manufacture of American blend cigarettes, moist snuff, chewing tobacco, and pipe tobacco. Liquorice provides tobacco products with a natural sweetness and a distinctive flavour that blends readily with the natural and imitation flavouring components employed in the tobacco industry. Liquorice can also be added to cigarette rolling papers. , the US Food and Drug Administration banned the use of any "characterizing flavors" other than menthol from cigarettes, but not other manufactured tobacco products.
Food and confectionery
Liquorice flavour is found in a wide variety of candies or sweets. In most of these candies, the taste is reinforced by aniseed oil so the actual content of liquorice is low.
In the Netherlands, liquorice confectionery (drop) is a common sweet sold in many forms. Mixing it with mint, menthol, aniseed, or laurel is common. It is also mixed with ammonium chloride (); salmiak liquorice in the Netherlands is known as ('salty liquorice'). Strong, salty sweets are also consumed in Nordic countries where liquorice flavoured alcohols are sold, particularly in Denmark and Finland.
Dried sticks of the liquorice root are a traditional confectionery in the Netherlands as were they once in Britain. They were sold simply as sticks of ('sweet wood') to chew on as a candy.
Pontefract in Yorkshire, England, is where liquorice mixed with sugar began to be used as a sweet in the contemporary way. Pontefract cakes were originally made there. In Cumbria, County Durham, Yorkshire and Lancashire, it is colloquially known as 'Spanish', supposedly because Spanish monks grew liquorice root at Rievaulx Abbey near Thirsk.
In Italy, Spain and France, liquorice is used in its natural form. The root of the plant is simply harvested, washed, dried, and chewed as a mouth freshener. Throughout Italy, unsweetened liquorice is consumed in the form of small black pieces made only from 100% pure liquorice extract. In Calabria, a liqueur is made from pure liquorice extract and in Reggio Emilia a soft drink called acqua d'orcio is made. In some parts of the Arab world, including Syria, Egypt, and Palestine, it is consumed as a cold beverage, especially in Ramadan.
In southeastern Turkey, such as in Diyarbakır, licorice root is traditionally made into a chilled beverage that is most commonly consumed in summer.
Research
Properties of glycyrrhizin are under preliminary research, such as for hepatitis C or topical treatment of psoriasis, but the low quality of studies prevents conclusions about efficacy and safety.
Traditional medicine
In traditional Chinese medicine, a related species G. uralensis (often translated as "liquorice") is known as (), and is believed to "harmonize" the ingredients in a formula. although there is no high-quality clinical research to indicate it is safe or effective for any medicinal purpose. The European Medical Agency added liquorice to their list of herbal medicine.
Fungicide
The essential oils inhibit the growth of Aspergillus flavus.
Adverse effects
Consumption levels
The United States Food and Drug Administration regards that foods containing liquorice and its derivatives (including glycyrrhizin) are generally recognized as safe for use as a food ingredient, if not consumed excessively. Other jurisdictions have suggested no more than of glycyrrhizin per day, the equivalent of about of liquorice confectionery. Although liquorice is considered safe as a food ingredient, glycyrrhizin can cause serious side effects if consumed in large amounts (above 0.2 mg per kg per day). One estimate is that a normal healthy person can consume of glycyrrhizic acid per day.
Because the composition of liquorice extracts in various products may exist in a broad range, there is not enough scientific information to determine that a specific level of intake is safe or unsafe.
Physiological effects
The effects of excessive liquorice consumption on lowering potassium levels in the blood and increasing blood pressure are a particular concern for people with hypertension (high blood pressure) or heart or kidney disease.
Some adverse effects of liquorice consumed in amounts of 50 to 200 g per day over four weeks appear to be caused by glycyrrhizic acid (75 to 540 mg per day glycyrrhetinic acid) causing increases in blood pressure. Consuming large amounts of liquorice during pregnancy has been associated with premature birth and health problems in the child.
Hyper-mineralocorticosteroid syndrome can occur when the body retains sodium, and loses potassium, altering biochemical and hormonal regulation. Some of these activities may include raised aldosterone levels, decline of the renin-angiotensin system and increased levels of the atrial natriuretic hormone in order to compensate the variations in homoeostasis.
Other adverse effects may include electrolyte imbalance, edema, increased blood pressure, weight gain, heart problems, and weakness. Symptoms depend on the severity of toxicity. Some other complaints include fatigue, shortness of breath, kidney failure, and paralysis.
Potential for toxicity
The major dose-limiting toxicities of liquorice are corticosteroid in nature, because of the inhibitory effect that its chief active constituents, glycyrrhizin and enoxolone, have on cortisol degradation, and include edema, hypokalaemia, weight gain or loss, and hypertension.
Gallery
| Biology and health sciences | Fabales | null |
249254 | https://en.wikipedia.org/wiki/Clique%20problem | Clique problem | In computer science, the clique problem is the computational problem of finding cliques (subsets of vertices, all adjacent to each other, also called complete subgraphs) in a graph. It has several different formulations depending on which cliques, and what information about the cliques, should be found. Common formulations of the clique problem include finding a maximum clique (a clique with the largest possible number of vertices), finding a maximum weight clique in a weighted graph, listing all maximal cliques (cliques that cannot be enlarged), and solving the decision problem of testing whether a graph contains a clique larger than a given size.
The clique problem arises in the following real-world setting. Consider a social network, where the graph's vertices represent people, and the graph's edges represent mutual acquaintance. Then a clique represents a subset of people who all know each other, and algorithms for finding cliques can be used to discover these groups of mutual friends. Along with its applications in social networks, the clique problem also has many applications in bioinformatics, and computational chemistry.
Most versions of the clique problem are hard. The clique decision problem is NP-complete (one of Karp's 21 NP-complete problems). The problem of finding the maximum clique is both fixed-parameter intractable and hard to approximate. And, listing all maximal cliques may require exponential time as there exist graphs with exponentially many maximal cliques. Therefore, much of the theory about the clique problem is devoted to identifying special types of graph that admit more efficient algorithms, or to establishing the computational difficulty of the general problem in various models of computation.
To find a maximum clique, one can systematically inspect all subsets, but this sort of brute-force search is too time-consuming to be practical for networks comprising more than a few dozen vertices.
Although no polynomial time algorithm is known for this problem, more efficient algorithms than the brute-force search are known. For instance, the Bron–Kerbosch algorithm can be used to list all maximal cliques in worst-case optimal time, and it is also possible to list them in polynomial time per clique.
History and applications
The study of complete subgraphs in mathematics predates the "clique" terminology. For instance, complete subgraphs make an early appearance in the mathematical literature in the graph-theoretic reformulation of Ramsey theory by . But the term "clique" and the problem of algorithmically listing cliques both come from the social sciences, where complete subgraphs are used to model social cliques, groups of people who all know each other. used graphs to model social networks, and adapted the social science terminology to graph theory. They were the first to call complete subgraphs "cliques". The first algorithm for solving the clique problem is that of , who were motivated by the sociological application. Social science researchers have also defined various other types of cliques and maximal cliques in social network, "cohesive subgroups" of people or actors in the network all of whom share one of several different kinds of connectivity relation. Many of these generalized notions of cliques can also be found by constructing an undirected graph whose edges represent related pairs of actors from the social network, and then applying an algorithm for the clique problem to this graph.
Since the work of Harary and Ross, many others have devised algorithms for various versions of the clique problem. In the 1970s, researchers began studying these algorithms from the point of view of worst-case analysis. See, for instance, , an early work on the worst-case complexity of the maximum clique problem. Also in the 1970s, beginning with the work of and , researchers began using the theory of NP-completeness and related intractability results to provide a mathematical explanation for the perceived difficulty of the clique problem. In the 1990s, a breakthrough series of papers beginning with showed that (assuming P ≠ NP) it is not even possible to approximate the problem accurately and efficiently.
Clique-finding algorithms have been used in chemistry, to find chemicals that match a target structure and to model molecular docking and the binding sites of chemical reactions. They can also be used to find similar structures within different molecules. In these applications, one forms a graph in which each vertex represents a matched pair of atoms, one from each of two molecules. Two vertices are connected by an edge if the matches that they represent are compatible with each other. Being compatible may mean, for instance, that the distances between the atoms within the two molecules are approximately equal, to within some given tolerance. A clique in this graph represents a set of matched pairs of atoms in which all the matches are compatible with each other. A special case of this method is the use of the modular product of graphs to reduce the problem of finding the maximum common induced subgraph of two graphs to the problem of finding a maximum clique in their product.
In automatic test pattern generation, finding cliques can help to bound the size of a test set. In bioinformatics, clique-finding algorithms have been used to infer evolutionary trees, predict protein structures, and find closely interacting clusters of proteins. Listing the cliques in a dependency graph is an important step in the analysis of certain random processes. In mathematics, Keller's conjecture on face-to-face tiling of hypercubes was disproved by , who used a clique-finding algorithm on an associated graph to find a counterexample.
Definitions
An undirected graph is formed by a finite set of vertices and a set of unordered pairs of vertices, which are called edges. By convention, in algorithm analysis, the number of vertices in the graph is denoted by and the number of edges is denoted by . A clique in a graph is a complete subgraph of . That is, it is a subset of the vertices such that every two vertices in are the two endpoints of an edge in . A maximal clique is a clique to which no more vertices can be added. For each vertex that is not part of a maximal clique, there must be another vertex that is in the clique and non-adjacent to , preventing from being added to the clique. A maximum clique is a clique that includes the largest possible number of vertices. The clique number is the number of vertices in a maximum clique of .
Several closely related clique-finding problems have been studied.
In the maximum clique problem, the input is an undirected graph, and the output is a maximum clique in the graph. If there are multiple maximum cliques, one of them may be chosen arbitrarily.
In the weighted maximum clique problem, the input is an undirected graph with weights on its vertices (or, less frequently, edges) and the output is a clique with maximum total weight. The maximum clique problem is the special case in which all weights are equal. As well as the problem of optimizing the sum of weights, other more complicated bicriterion optimization problems have also been studied.
In the maximal clique listing problem, the input is an undirected graph, and the output is a list of all its maximal cliques. The maximum clique problem may be solved using as a subroutine an algorithm for the maximal clique listing problem, because the maximum clique must be included among all the maximal cliques.
In the -clique problem, the input is an undirected graph and a number . The output is a clique with vertices, if one exists, or a special value indicating that there is no -clique otherwise. In some variations of this problem, the output should list all cliques of size .
In the clique decision problem, the input is an undirected graph and a number , and the output is a Boolean value: true if the graph contains a -clique, and false otherwise.
The first four of these problems are all important in practical applications. The clique decision problem is not of practical importance; it is formulated in this way in order to apply the theory of NP-completeness to clique-finding problems.
The clique problem and the independent set problem are complementary: a clique in is an independent set in the complement graph of and vice versa. Therefore, many computational results may be applied equally well to either problem, and some research papers do not clearly distinguish between the two problems. However, the two problems have different properties when applied to restricted families of graphs. For instance, the clique problem may be solved in polynomial time for planar graphs while the independent set problem remains NP-hard on planar graphs.
Algorithms
Finding a single maximal clique
A maximal clique, sometimes called inclusion-maximal, is a clique that is not included in a larger clique. Therefore, every clique is contained in a maximal clique. Maximal cliques can be very small. A graph may contain a non-maximal clique with many vertices and a separate clique of size 2 which is maximal. While a maximum (i.e., largest) clique is necessarily maximal, the converse does not hold. There are some types of graphs in which every maximal clique is maximum; these are the complements of the well-covered graphs, in which every maximal independent set is maximum. However, other graphs have maximal cliques that are not maximum.
A single maximal clique can be found by a straightforward greedy algorithm. Starting with an arbitrary clique (for instance, any single vertex or even the empty set), grow the current clique one vertex at a time by looping through the graph's remaining vertices. For each vertex that this loop examines, add to the clique if it is adjacent to every vertex that is already in the clique, and discard otherwise. This algorithm runs in linear time. Because of the ease of finding maximal cliques, and their potential small size, more attention has been given to the much harder algorithmic problem of finding a maximum or otherwise large clique. However, some research in parallel algorithms has studied the problem of finding a maximal clique. In particular, the problem of finding the lexicographically first maximal clique (the one found by the algorithm above) has been shown to be complete for the class of polynomial-time functions. This result implies that the problem is unlikely to be solvable within the parallel complexity class NC.
Cliques of fixed size
One can test whether a graph contains a -vertex clique, and find any such clique that it contains, using a brute force algorithm. This algorithm examines each subgraph with vertices and checks to see whether it forms a clique. It takes time , as expressed using big O notation.
This is because there are subgraphs to check, each of which has edges whose presence in needs to be checked. Thus, the problem may be solved in polynomial time whenever is a fixed constant. However, when does not have a fixed value, but instead may vary as part of the input to the problem, the time is exponential.
The simplest nontrivial case of the clique-finding problem is finding a triangle in a graph, or equivalently determining whether the graph is triangle-free.
In a graph with edges, there may be at most triangles (using big theta notation to indicate that this bound is tight). The worst case for this formula occurs when is itself a clique. Therefore, algorithms for listing all triangles must take at least time in the worst case (using big omega notation), and algorithms are known that match this time bound. For instance, describe an algorithm that sorts the vertices in order from highest degree to lowest and then iterates through each vertex in the sorted list, looking for triangles that include and do not include any previous vertex in the list. To do so the algorithm marks all neighbors of , searches through all edges incident to a neighbor of outputting a triangle for every edge that has two marked endpoints, and then removes the marks and deletes from the graph. As the authors show, the time for this algorithm is proportional to the arboricity of the graph (denoted ) multiplied by the number of edges, which is . Since the arboricity is at most , this algorithm runs in time . More generally, all -vertex cliques can be listed by a similar algorithm that takes time proportional to the number of edges multiplied by the arboricity to the power . For graphs of constant arboricity, such as planar graphs (or in general graphs from any non-trivial minor-closed graph family), this algorithm takes time, which is optimal since it is linear in the size of the input.
If one desires only a single triangle, or an assurance that the graph is triangle-free, faster algorithms are possible. As observe, the graph contains a triangle if and only if its adjacency matrix and the square of the adjacency matrix contain nonzero entries in the same cell. Therefore, fast matrix multiplication techniques can be applied to find triangles in time . used fast matrix multiplication to improve the algorithm for finding triangles to . These algorithms based on fast matrix multiplication have also been extended to problems of finding -cliques for larger values of .
Listing all maximal cliques
By a result of , every -vertex graph has at most maximal cliques. They can be listed by the Bron–Kerbosch algorithm, a recursive backtracking procedure of . The main recursive subroutine of this procedure has three arguments: a partially constructed (non-maximal) clique, a set of candidate vertices that could be added to the clique, and another set of vertices that should not be added (because doing so would lead to a clique that has already been found). The algorithm tries adding the candidate vertices one by one to the partial clique, making a recursive call for each one. After trying each of these vertices, it moves it to the set of vertices that should not be added again. Variants of this algorithm can be shown to have worst-case running time , matching the number of cliques that might need to be listed. Therefore, this provides a worst-case-optimal solution to the problem of listing all maximal cliques. Further, the Bron–Kerbosch algorithm has been widely reported as being faster in practice than its alternatives.
However, when the number of cliques is significantly smaller than its worst case, other algorithms might be preferable. As showed, it is also possible to list all maximal cliques in a graph in an amount of time that is polynomial per generated clique. An algorithm such as theirs in which the running time depends on the output size is known as an output-sensitive algorithm. Their algorithm is based on the following two observations, relating the maximal cliques of the given graph to the maximal cliques of a graph formed by removing an arbitrary vertex from :
For every maximal clique of , either continues to form a maximal clique in , or forms a maximal clique in . Therefore, has at least as many maximal cliques as does.
Each maximal clique in that does not contain is a maximal clique in , and each maximal clique in that does contain can be formed from a maximal clique in by adding and removing the non-neighbors of from .
Using these observations they can generate all maximal cliques in by a recursive algorithm that chooses a vertex arbitrarily and then, for each maximal clique in , outputs both and the clique formed by adding to and removing the non-neighbors of . However, some cliques of may be generated in this way from more than one parent clique of , so they eliminate duplicates by outputting a clique in only when its parent in is lexicographically maximum among all possible parent cliques. On the basis of this principle, they show that all maximal cliques in may be generated in time per clique, where is the number of edges in and is the number of vertices. improve this to per clique, where is the arboricity of the given graph. provide an alternative output-sensitive algorithm based on fast matrix multiplication. show that it is even possible to list all maximal cliques in lexicographic order with polynomial delay per clique. However, the choice of ordering is important for the efficiency of this algorithm: for the reverse of this order,
there is no polynomial-delay algorithm unless P = NP.
On the basis of this result, it is possible to list all maximal cliques in polynomial time, for families of graphs in which the number of cliques is polynomially bounded. These families include chordal graphs, complete graphs, triangle-free graphs, interval graphs, graphs of bounded boxicity, and planar graphs. In particular, the planar graphs have cliques, of at most constant size, that can be listed in linear time. The same is true for any family of graphs that is both sparse (having a number of edges at most a constant times the number of vertices) and closed under the operation of taking subgraphs.
Finding maximum cliques in arbitrary graphs
It is possible to find the maximum clique, or the clique number, of an arbitrary n-vertex graph in time by using one of the algorithms described above to list all maximal cliques in the graph and returning the largest one. However, for this variant of the clique problem better worst-case time bounds are possible. The algorithm of solves this problem in time . It is a recursive backtracking scheme similar to that of the Bron–Kerbosch algorithm, but is able to eliminate some recursive calls when it can be shown that the cliques found within the call will be suboptimal. improved the time to , and improved it to time, at the expense of greater space usage. Robson's algorithm combines a similar backtracking scheme (with a more complicated case analysis) and a dynamic programming technique in which the optimal solution is precomputed for all small connected subgraphs of the complement graph. These partial solutions are used to shortcut the backtracking recursion. The fastest algorithm known today is a refined version of this method by which runs in time .
There has also been extensive research on heuristic algorithms for solving maximum clique problems without worst-case runtime guarantees, based on methods including branch and bound, local search, greedy algorithms, and constraint programming. Non-standard computing methodologies that have been suggested for finding cliques include DNA computing and adiabatic quantum computation. The maximum clique problem was the subject of an implementation challenge sponsored by DIMACS in 1992–1993, and a collection of graphs used as benchmarks for the challenge, which is publicly available.
Special classes of graphs
Planar graphs, and other families of sparse graphs, have been discussed above: they have linearly many maximal cliques, of bounded size, that can be listed in linear time. In particular, for planar graphs, any clique can have at most four vertices, by Kuratowski's theorem.
Perfect graphs are defined by the properties that their clique number equals their chromatic number, and that this equality holds also in each of their induced subgraphs. For perfect graphs, it is possible to find a maximum clique in polynomial time, using an algorithm based on semidefinite programming.
However, this method is complex and non-combinatorial, and specialized clique-finding algorithms have been developed for many subclasses of perfect graphs. In the complement graphs of bipartite graphs, Kőnig's theorem allows the maximum clique problem to be solved using techniques for matching. In another class of perfect graphs, the permutation graphs, a maximum clique is a longest decreasing subsequence of the permutation defining the graph and can be found using known algorithms for the longest decreasing subsequence problem. Conversely, every instance of the longest decreasing subsequence problem can be described equivalently as a problem of finding a maximum clique in a permutation graph. provide an alternative quadratic-time algorithm for maximum cliques in comparability graphs, a broader class of perfect graphs that includes the permutation graphs as a special case. In chordal graphs, the maximal cliques can be found by listing the vertices in an elimination ordering, and checking the clique neighborhoods of each vertex in this ordering.
In some cases, these algorithms can be extended to other, non-perfect, classes of graphs as well. For instance, in a circle graph, the neighborhood of each vertex is a permutation graph, so a maximum clique in a circle graph can be found by applying the permutation graph algorithm to each neighborhood. Similarly, in a unit disk graph (with a known geometric representation), there is a polynomial time algorithm for maximum cliques based on applying the algorithm for complements of bipartite graphs to shared neighborhoods of pairs of vertices.
The algorithmic problem of finding a maximum clique in a random graph drawn from the Erdős–Rényi model (in which each edge appears with probability , independently from the other edges) was suggested by . Because the maximum clique in a random graph has logarithmic size with high probability, it
can be found by a brute force search in expected time . This is a quasi-polynomial time bound. Although the clique number of such graphs is usually very close to , simple greedy algorithms as well as more sophisticated randomized approximation techniques only find cliques with size , half as big. The number of maximal cliques in such graphs is with high probability exponential in , which prevents methods that list all maximal cliques from running in polynomial time. Because of the difficulty of this problem, several authors have investigated the planted clique problem, the clique problem on random graphs that have been augmented by adding large cliques. While spectral methods and semidefinite programming can detect hidden cliques of size , no polynomial-time algorithms are currently known to detect those of size (expressed using little-o notation).
Approximation algorithms
Several authors have considered approximation algorithms that attempt to find a clique or independent set that, although not maximum, has size as close to the maximum as can be found in polynomial time. Although much of this work has focused on independent sets in sparse graphs, a case that does not make sense for the complementary clique problem, there has also been work on approximation algorithms that do not use such sparsity assumptions.
describes a polynomial time algorithm that finds a clique of size in any graph that has clique number for any constant . By using this algorithm when the clique number of a given input graph is between and , switching to a different algorithm of for graphs with higher clique numbers, and choosing a two-vertex clique if both algorithms fail to find anything, Feige provides an approximation algorithm that finds a clique with a number of vertices within a factor of of the maximum. Although the approximation ratio of this algorithm is weak, it is the best known to date. The results on hardness of approximation described below suggest that there can be no approximation algorithm with an approximation ratio significantly less than linear.
Lower bounds
NP-completeness
The clique decision problem is NP-complete. It was one of Richard Karp's original 21 problems shown NP-complete in his 1972 paper "Reducibility Among Combinatorial Problems". This problem was also mentioned in Stephen Cook's paper introducing the theory of NP-complete problems. Because of the hardness of the decision problem, the problem of finding a maximum clique is also NP-hard. If one could solve it, one could also solve the decision problem, by comparing the size of the maximum clique to the size parameter given as input in the decision problem.
Karp's NP-completeness proof is a many-one reduction from the Boolean satisfiability problem. It describes how to translate Boolean formulas in conjunctive normal form (CNF) into equivalent instances of the maximum clique problem. Satisfiability, in turn, was proved NP-complete in the Cook–Levin theorem. From a given CNF formula, Karp forms a graph that has a vertex for every pair , where is a variable or its negation and is a clause in the formula that contains . Two of these vertices are connected by an edge if they represent compatible variable assignments for different clauses. That is, there is an edge from to whenever and and are not each other's negations. If denotes the number of clauses in the CNF formula, then the -vertex cliques in this graph represent consistent ways of assigning truth values to some of its variables in order to satisfy the formula. Therefore, the formula is satisfiable if and only if a -vertex clique exists.
Some NP-complete problems (such as the travelling salesman problem in planar graphs) may be solved in time that is exponential in a sublinear function of the input size parameter , significantly faster than a brute-force search. However, it is unlikely that such a subexponential time bound is possible for the clique problem in arbitrary graphs, as it would imply similarly subexponential bounds for many other standard NP-complete problems.
Circuit complexity
The computational difficulty of the clique problem has led it to be used to prove several lower bounds in circuit complexity. The existence of a clique of a given size is a monotone graph property, meaning that, if a clique exists in a given graph, it will exist in any supergraph. Because this property is monotone, there must exist a monotone circuit, using only and gates and or gates, to solve the clique decision problem for a given fixed clique size. However, the size of these circuits can be proven to be a super-polynomial function of the number of vertices and the clique size, exponential in the cube root of the number of vertices. Even if a small number of NOT gates are allowed, the complexity remains superpolynomial. Additionally, the depth of a monotone circuit for the clique problem using gates of bounded fan-in must be at least a polynomial in the clique size.
Decision tree complexity
The (deterministic) decision tree complexity of determining a graph property is the number of questions of the form "Is there an edge between vertex and vertex ?" that have to be answered in the worst case to determine whether a graph has a particular property. That is, it is the minimum height of a Boolean decision tree for the problem. There are possible questions to be asked. Therefore, any graph property can be determined with at most questions. It is also possible to define random and quantum decision tree complexity of a property, the expected number of questions (for a worst case input) that a randomized or quantum algorithm needs to have answered in order to correctly determine whether the given graph has the property.
Because the property of containing a clique is monotone, it is covered by the Aanderaa–Karp–Rosenberg conjecture, which states that the deterministic decision tree complexity of determining any non-trivial monotone graph property is exactly . For arbitrary monotone graph properties, this conjecture remains unproven. However, for deterministic decision trees, and for any in the range , the property of containing a -clique was shown to have decision tree complexity exactly by . Deterministic decision trees also require exponential size to detect cliques, or large polynomial size to detect cliques of bounded size.
The Aanderaa–Karp–Rosenberg conjecture also states that the randomized decision tree complexity of non-trivial monotone functions is . The conjecture again remains unproven, but has been resolved for the property of containing a clique for . This property is known to have randomized decision tree complexity . For quantum decision trees, the best known lower bound is , but no matching algorithm is known for the case of .
Fixed-parameter intractability
Parameterized complexity is the complexity-theoretic study of problems that are naturally equipped with a small integer parameter and for which the problem becomes more difficult as increases, such as finding -cliques in graphs. A problem is said to be fixed-parameter tractable if there is an algorithm for solving it on inputs of size , and a function , such that the algorithm runs in time . That is, it is fixed-parameter tractable if it can be solved in polynomial time for any fixed value of and moreover if the exponent of the polynomial does not depend on .
For finding -vertex cliques, the brute force search algorithm has running time . Because the exponent of depends on , this algorithm is not fixed-parameter tractable.
Although it can be improved by fast matrix multiplication, the running time still has an exponent that is linear in . Thus, although the running time of known algorithms for the clique problem is polynomial for any fixed , these algorithms do not suffice for fixed-parameter tractability. defined a hierarchy of parametrized problems, the W hierarchy, that they conjectured did not have fixed-parameter tractable algorithms. They proved that independent set (or, equivalently, clique) is hard for the first level of this hierarchy, W[1]. Thus, according to their conjecture, clique has no fixed-parameter tractable algorithm. Moreover, this result provides the basis for proofs of W[1]-hardness of many other problems, and thus serves as an analogue of the Cook–Levin theorem for parameterized complexity.
showed that finding -vertex cliques cannot be done in time unless the exponential time hypothesis fails. Again, this provides evidence that no fixed-parameter tractable algorithm is possible.
Although the problems of listing maximal cliques or finding maximum cliques are unlikely to be fixed-parameter tractable with the parameter , they may be fixed-parameter tractable for other parameters of instance complexity. For instance, both problems are known to be fixed-parameter tractable when parametrized by the degeneracy of the input graph.
Hardness of approximation
Weak results hinting that the clique problem might be hard to approximate have been known for a long time. observed that, because the clique number takes on small integer values and is NP-hard to compute, it cannot have a fully polynomial-time approximation scheme, unless P = NP. If too accurate an approximation were available, rounding its value to an integer would give the exact clique number. However, little more was known until the early 1990s, when several authors began to make connections between the approximation of maximum cliques and probabilistically checkable proofs. They used these connections to prove hardness of approximation results for the maximum clique problem.
After many improvements to these results it is now known that, for every real number , there can be no polynomial time algorithm that approximates the maximum clique to within a factor better than , unless P = NP.
The rough idea of these inapproximability results is to form a graph that represents a probabilistically checkable proof system for an NP-complete problem such as the Boolean satisfiability problem. In a probabilistically checkable proof system, a proof is represented as a sequence of bits. An instance of the satisfiability problem should have a valid proof if and only if it is satisfiable. The proof is checked by an algorithm that, after a polynomial-time computation on the input to the satisfiability problem, chooses to examine a small number of randomly chosen positions of the proof string. Depending on what values are found at that sample of bits, the checker will either accept or reject the proof, without looking at the rest of the bits. False negatives are not allowed: a valid proof must always be accepted. However, an invalid proof may sometimes mistakenly be accepted. For every invalid proof, the probability that the checker will accept it must be low.
To transform a probabilistically checkable proof system of this type into a clique problem, one forms a graph with a vertex for each possible accepting run of the proof checker. That is, a vertex is defined by one of the possible random choices of sets of positions to examine, and by bit values for those positions that would cause the checker to accept the proof. It can be represented by a partial word with a 0 or 1 at each examined position and a wildcard character at each remaining position. Two vertices are adjacent, in this graph, if the corresponding two accepting runs see the same bit values at every position they both examine. Each (valid or invalid) proof string corresponds to a clique, the set of accepting runs that see that proof string, and all maximal cliques arise in this way. One of these cliques is large if and only if it corresponds to a proof string that many proof checkers accept. If the original satisfiability instance is satisfiable, it will have a valid proof string, one that is accepted by all runs of the checker, and this string will correspond to a large clique in the graph. However, if the original instance is not satisfiable, then all proof strings are invalid, each proof string has only a small number of checker runs that mistakenly accept it, and all cliques are small. Therefore, if one could distinguish in polynomial time between graphs that have large cliques and graphs in which all cliques are small, or if one could accurately approximate the clique problem, then applying this approximation to the graphs generated from satisfiability instances would allow satisfiable instances to be distinguished from unsatisfiable instances. However, this is not possible unless P = NP.
| Mathematics | Graph theory | null |
249357 | https://en.wikipedia.org/wiki/Sleepwalking | Sleepwalking | Sleepwalking, also known as somnambulism or noctambulism, is a phenomenon of combined sleep and wakefulness. It is classified as a sleep disorder belonging to the parasomnia family. It occurs during the slow wave stage of sleep, in a state of low consciousness, with performance of activities that are usually performed during a state of full consciousness. These activities can be as benign as talking, sitting up in bed, walking to a bathroom, consuming food, and cleaning, or as hazardous as cooking, driving a motor vehicle, violent gestures and grabbing at hallucinated objects.
Although sleepwalking cases generally consist of simple, repeated behaviors, there are occasionally reports of people performing complex behaviors while asleep, although their legitimacy is often disputed. Sleepwalkers often have little or no memory of the incident, as their consciousness has altered into a state in which memories are difficult to recall. Although their eyes are open, their expression is dim and glazed over. This may last from 30 seconds to 30 minutes.
Sleepwalking occurs during slow-wave sleep (N3) of non-rapid eye movement sleep (NREM sleep) cycles. It typically occurs within the first third of the night when slow-wave sleep is most prominent. Usually, it will occur once in a night, if at all.
Signs and symptoms
Sleepwalking is characterized by:
partial arousal during non-rapid eye movement (NREM) sleep, typically during the first third of the night
dreamy content that may or may not be recalled when awake
dream-congruent motor behavior that may be simple or complex
impaired perception of the environment
impaired judgement, planning and problem-solving.
Despite how it is portrayed in many cultures (eyes closed and walking with arms outstretched), the sleepwalker's eyes are open but may appear as a glassy-eyed stare or blank expression and pupils are dilated. They are often disoriented, consequent to awakening: the sleepwalker may be confused and perplexed, and might not know why or how they got out of bed; however, the disorientation will fade within minutes. They may talk while sleepwalking, but the talk typically does not make sense to the observer. There are varying degrees of amnesia associated with sleepwalking, ranging from no memory at all, vague memories or a narrative.
Associated disorders
Most studies look at sleep disorders in adults but children can also be affected. In the ten percent of the population that experience sleep-related disorders, children are mainly affected due to their youthful brains. A study conducted in Australia, looked at sleepwalking and its association with sleep behaviors in children. It was found that sleepwalking could be associated with children's bedtime routines. Those who have behavioral problems are more likely to develop a sleep disorder and should be assessed. The relationship between sleepwalking and the behavioral and emotional problems are more associated than their bedtime routines. This may very well be because sleep related disorders and sleepwalking happen simultaneously; one cannot exist without the other.
In the study "Sleepwalking and Sleep Terrors in Prepubertal Children" it was found that, if a child had another sleep disordersuch as restless leg syndrome (RLS) or sleep-disorder breathing (SDB)there was a greater chance of sleepwalking. The study found that children with chronic parasomnias may often also present SDB or, to a lesser extent, RLS. Furthermore, the disappearance of the parasomnias after the treatment of the SDB or RLS periodic limb movement syndrome suggests that the latter may trigger the former. The high frequency of SDB in family members of children with parasomnia provided additional evidence that SDB may manifest as parasomnias in children. Children with parasomnias are not systematically monitored during sleep, although past studies have suggested that patients with sleep terrors or sleepwalking have an elevated level of brief EEG arousals. When children receive polysomnographies, discrete patterns (e.g., nasal flow limitation, abnormal respiratory effort, bursts of high or slow EEG frequencies) should be sought; apneas are rarely found in children. Children's respiration during sleep should be monitored with nasal cannula or pressure transducer system or esophageal manometry, which are more sensitive than the thermistors or thermocouples currently used in many laboratories. The clear, prompt improvement of severe parasomnia in children who are treated for SDB, as defined here, provides important evidence that subtle SDB can have substantial health-related significance. Also noteworthy is the report of familial presence of parasomnia. Studies of twin cohorts and families with sleep terror and sleepwalking suggest genetic involvement of parasomnias. RLS and SDB have been shown to have familial recurrence. RLS has been shown to have genetic involvement.
Sleepwalking may also accompany the related phenomenon of night terrors, especially in children. In the midst of a night terror, the affected person may wander in a distressed state while still asleep, and examples of sufferers attempting to run or aggressively defend themselves during these incidents have been reported in medical literature.
In some cases, sleepwalking in adults may be a symptom of a psychological disorder. One study suggests higher levels of dissociation in adult sleepwalkers, since test subjects scored unusually high on the hysteria portion of the "Crown-Crisp Experiential Index". Another suggested that "A higher incidence [of sleepwalking events] has been reported in patients with schizophrenia, hysteria and anxiety neuroses". Also, patients with migraine headaches or Tourette syndrome are 4–6 times more likely to sleepwalk.
Consequences
During the amnesic state sleepwalkers are in, many things can happen without their recollection. One thing that can happen is a sleep disorder called sexomnia, where an individual can engage in sexual behaviors with oneself or others. Its occurrence is rare, but can happen during sleepwalking. Sleep-related eating disorder, in which sleepwalkers eat involuntarily, can also happen. The events can include eating/drinking regular foods or odd combinations of food. Insomnia and daytime sleepiness can also occur. Most sleepwalkers get injuries at some point during sleepwalking, often minor injuries such as cuts or bruises. In rare occasions, however, sleepwalkers have fractured bones and died as the result of a fall. Sleepwalkers may also face embarrassment of being found naked in public.
Causes
The cause of sleepwalking is unknown. A number of, as yet unproven, hypotheses are suggested for why it might occur, including: delay in the maturity of the central nervous system, increased slow wave sleep, sleep deprivation, fever, and excessive tiredness. There may be a genetic component to sleepwalking. One study found that sleepwalking occurred in 45% of children who have one parent who sleepwalked, and in 60% of children if both parents sleepwalked. Thus, heritable factors may predispose an individual to sleepwalking, but expression of the behavior may also be influenced by environmental factors. Genetic studies using common fruit flies as experimental models reveal a link between night sleep and brain development mediated by evolutionary conserved transcription factors such as AP-2 Sleepwalking may be inherited as an autosomal dominant disorder with reduced penetrance. Genome-wide multipoint parametric linkage analysis for sleepwalking revealed a maximum logarithm of the odds score of 3.14 at chromosome 20q12-q13.12 between 55.6 and 61.4 cM.
Sleepwalking has been hypothesized to be linked to the neurotransmitter serotonin, which also appears to be metabolized differently in migraine patients and people with Tourette syndrome, both populations being four to nine times more likely to experience an episode of sleepwalking. Hormonal fluctuations have been found to contribute to sleepwalking episodes in women, with the likeliness to sleepwalk being higher before the onset of menstruation. It also appears that hormonal changes during pregnancy decrease the likelihood of engaging in sleepwalking.
Medications, primarily in four classes—benzodiazepine receptor agonists and other GABA modulators, antidepressants and other serotonergic agents, antipsychotics, and β-blockers—have been associated with sleepwalking. The best evidence of medications causing sleepwalking is for zolpidem and sodium oxybate; all other reports are based on associations noted in case reports.
A number of conditions, such as Parkinson's disease, are thought to trigger sleepwalking in people without a previous history of sleepwalking.
Diagnosis
Polysomnography is the only accurate assessment of a sleepwalking episode. Because this is costly and sleepwalking episodes are usually infrequent, other measures commonly used include self-, parent-, or partner-report. Three common diagnostic systems that are generally used for sleepwalking disorders are International Classification of Diseases (ICD-10), the International Classification of Sleep Disorders (ICSD-3), and the Diagnostic and Statistical Manual.
The Diagnostic and Statistical Manual defines two subcategories of sleepwalking, although sleepwalking does not need to involve either behaviours:
sleepwalking with sleep-related eating.
sleepwalking with sleep-related sexual behavior (sexsomnia).
Sleep eating involves consuming food while asleep. These sleep eating disorders are more often than not induced for stress related reasons. Another major cause of this sleep eating subtype of sleepwalking is sleep medication, such as Ambien for example (Mayo Clinic). There are a few others, but Ambien is a more widely used sleep aid. Because many sleep eaters prepare the food they consume, there are risks involving burns and such with ovens and other appliances. As expected, weight gain is also a common outcome of this disorder, because food that is frequently consumed contains high carbohydrates. As with sleepwalking, there are ways that sleep eating disorders can be maintained. There are some medications that calm the sleeper so they can get longer and better-quality rest, but activities such as yoga can also be introduced to reduce the stress and anxiety causing the action.
Differential diagnoses
Sleepwalking should not be confused with alcohol- or drug-induced blackouts, which can result in amnesia for events similar to sleepwalking. During an alcohol-induced blackout (drug-related amnesia), a person is able to actively engage and respond to their environment (e.g. having conversations or driving a vehicle), however the brain does not create memories for the events. Alcohol-induced blackouts can occur with blood alcohol levels higher than 0.06 g/dl. A systematic review of the literature found that approximately 50% of drinkers have experienced memory loss during a drinking episode and have had associated negative consequences similar to sleepwalkers, including injury and death.
Other differential diagnoses include rapid eye movement sleep behavior disorder, confusional arousals, and night terrors.
Assessment
An assessment of sleepwalking via polysomnography poses the problem that sleepwalking is less likely to occur in the sleep laboratory, and if an episode occurs, it is usually less complex than what the patient experiences at home. Therefore, the diagnosis can often be made by assessment of sleep history, time-course and content of the sleep related behaviors. Sometimes, home videos can provide additional information and should be considered in the diagnostic process.
Some features that should always be assessed include:
Age of onset
When the episode occurs during the sleep period
How often these episodes occur (frequency) and how long they last (duration)
Description of the episode, including behavior, emotions, and thoughts during and after the event
How responsive the patient is to external stimuli during the episode
How conscious or aware the patient is, when awakened from an episode
If the episode is remembered afterwards
The triggers or precipitating factors
Sleep–wake pattern and sleep environment
Daytime sleepiness
Other sleep disorders that might be present
Family history for NREM parasomnias and other sleep disorders
Medical, psychiatric, and neurological history
Medication and substance use history
The assessment should rule out differential diagnoses.
Treatment
There have been no clinical trials to show that any psychological or pharmacological intervention is effective in preventing sleepwalking episodes. Despite this, a wide range of treatments have been used with sleepwalkers. Psychological interventions have included psychoanalysis, hypnosis, scheduled or anticipatory waking, assertion training, relaxation training, managing aggressive feelings, sleep hygiene, classical conditioning (including electric shock), and play therapy. Pharmacological treatments have included tricyclic antidepressants (imipramine), an anticholinergic (biperiden), antiepileptics (carbamazepine, valproate), an antipsychotic (quetiapine), benzodiazepines (clonazepam, diazepam, flurazepam and triazolam), melatonin, a selective serotonin reuptake inhibitor (paroxetine), a barbiturate (sodium amytal) and herbs.
There is no evidence to show that waking sleepwalkers is harmful or not, though the sleepwalker is likely to be disoriented if awakened.
Unlike other sleep disorders, sleepwalking is not associated with daytime behavioral or emotional problems. This may be because the sleepwalker's sleep is not disturbed—unless they are woken, they are still in a sleep state while sleepwalking.
Maintaining the safety of the sleepwalker and others and seeking treatment for other sleep problems is recommended. Reassurance is recommended if sleepwalking is not causing any problems. However, if it causes distress or there is risk of harm, hypnosis and scheduled waking are recommended as treatments.
Safety planning
For those whose sleepwalking episodes are hazardous, a door alarm may offer a measure of protection. There are various kinds of door alarms that can attach to a bedroom door and when the door is opened, the alarm sounds. The intention is that the sound will fully awaken the person and interrupt the sleepwalking episode, or if the sleepwalker lives with others, the sound will prompt them to check on the person.
Sleepwalkers should aim to have their bedrooms on the ground floor of a home, apartment, dorm, hotel, etc.
Sleepwalkers should not have easily accessible weapons (loaded guns, knives) in the bedroom or any room of the house for that matter. If there are weapons, they should be locked away with keys secluded from the sleepwalker.
For partners of sleepwalkers who are violent or disturb their sleep, sleeping in another room may lead to better sleep quality and quantity.
Epidemiology
The lifetime prevalence of sleepwalking is estimated to be 4.6–10.3%. A meta-analysis of 51 studies, that included more than 100,000 children and adults, found that sleepwalking is more common in children with an estimated 5%, compared with 1.5% of adults, sleepwalking at least once in the previous 12 months. The rate of sleepwalking has not been found to vary across ages during childhood.
History
Sleepwalking has attracted a sense of mystery, but was not seriously investigated and diagnosed until the 19th century. The German chemist and parapsychologist Baron Karl Ludwig von Reichenbach (1788–1869) made extensive studies of sleepwalkers and used his discoveries to formulate his theory of the Odic force.
Sleepwalking was initially thought to be a dreamer acting out a dream. For example, in one study published by the Society for Science & the Public in 1954, this was the conclusion: "Repression of hostile feelings against the father caused the patients to react by acting out in a dream world with sleepwalking, the distorted fantasies they had about all authoritarian figures, such as fathers, officers and stern superiors." This same group published an article twelve years later with a new conclusion: "Sleepwalking, contrary to most belief, apparently has little to do with dreaming. In fact, it occurs when the sleeper is enjoying his most oblivious, deepest sleep—a stage in which dreams are not usually reported." More recent research has discovered that sleepwalking is actually a disorder of NREM (non-rapid eye movement) arousal. Acting out a dream is the basis for a REM (rapid eye movement) sleep disorder called REM Behavior Disorder (or REM Sleep Behavior Disorder). More accurate data about sleep is due to the invention of technologies, such as the electroencephalogram (EEG) by Hans Berger in 1924 and BEAM by Frank Duffy in the early 1980s.
In 1907, Sigmund Freud spoke about sleepwalking to the Vienna Psychoanalytic Society (Nunberg and Federn). He believed that sleepwalking was connected to fulfilling sexual wishes and was surprised that a person could move without interrupting their dream. At that time, Freud suggested that the essence of this phenomenon was the desire to go to sleep in the same area as the individual had slept in childhood. Ten years later, he speculated about somnambulism in the article "A Metapsychological Supplement to the Theory of Dreams" (1916–17 [1915]). In this essay, he clarified and expanded his hypothetical ideas on dreams. He described the dream as a fragile equilibrium that is destabilized by the repressed unconscious impulses of the unconscious system, which does not obey the wishes of the ego. Certain preconscious daytime thoughts can be resistant and these can retain a part of their cathexis as well. Unconscious impulses and day residues can come together and result in a conflict. Freud then wondered about the outcome of this wishful impulse: an unconscious instinctual demand that becomes a dream wish in the preconscious. Freud stated that this unconscious impulse could be expressed as mobility during sleep. This would be what is observed in somnambulism, though what actually makes it possible remains unknown.
As of 2002, sleepwalking has not been detected in non-human primates.
It is unclear whether it simply has not been observed yet, or whether sleepwalking is a uniquely human phenomenon.
Culture
Opera
Vincenzo Bellini's 1831 Italian opera semiseria, La sonnambula, the plot of which is centered on the question of the innocence of the betrothed and soon-to-be married Amina, who, upon having been discovered in the bedchamber of a stranger, and despite the assurances of that stranger that Amina was entirely innocent, has been rejected by her enraged fiancé, Elvino—who, then, decides to marry another. In fact, when stressed, Amina was susceptible to somnambulism; and had come to be in the stranger's bedchamber by sleep-walking along a high parapet (in full view of the opera's audience). Elvino, who later observes the (exhausted by all the fuss) Amina, sleep-walking across a very high, very unstable, and very rickety bridge at the local mill, realizes his mistake, abandons his plans of marriage to the other woman, and re-unites with Amina.
Jenny Lind and James Braid
In August 1847, the famous soprano Jenny Lind visited Manchester, and gave two performances as Amina. The outstanding difference between Lind and her contemporaries was that, "whilst the beauty of her voice was far greater than any other in living memory (thus, the Swedish Nightingale), what really set her apart was her outstanding ability to act"; and, moreover, in performing as Amina, rather than walking along a wide and well-protected walkway (as the others did), she routinely acrobatically balanced her way along narrow planks.
While she was in Manchester—on the basis that, at the time, many characterized "hypnotism" as "artificial somnambulism", and that, from a rather different perspective, her stage performance could also be described as one of "artificial" (rather than spontaneous) somnambulism—her friends arranged for her to visit the local surgeon James Braid, who had discovered hypnotism in 1841:
Drama
The sleepwalking scene (Act V Scene 1) from William Shakespeare's tragic play Macbeth (1606) is one of the most famous scenes in all of literature.
In Walley Chamberlain Oulton's two act farce The Sleep-Walker; or, Which is the Lady (1812), "Somno", a histrionic failed-actor-turned-manservant relives his wished-for roles when sleepwalking.
Literature
In Bram Stoker's novel Dracula, the character Lucy Westenra is described as a sleepwalker. It is while she is sleepwalking that Count Dracula lures and attacks her.
Sleepwalking as a legal defense
As sleepwalking behaviours occur without volition, sleepwalking can be used as a legal defense, as a form of legal automatism. An individual can be accused of non-insane or insane automatism. The first is used as a defense for temporary insanity or involuntary conduct, resulting in acquittal. The latter results in a "special verdict of not guilty by reason of insanity." This verdict of insanity can result in a court order to attend a mental institution.
In the 1963 case Bratty v A-G for Northern Ireland, Lord Morris stated, "Each set of facts must require a careful examination of its own circumstances, but if by way of taking an illustration it were considered possible for a person to walk in his sleep and to commit a violent crime while genuinely unconscious, then such a person would not be criminally liable for that act." While the veracity of the cases are disputed, there have been acts of homicide where the prime suspect may have committed the act while sleepwalking.
Alternative explanations to homicidal or violent sleepwalking include malingering, drug-induced amnesia, and other disorders in which sleep-related violence may occur, such as REM Sleep Behavior Disorder, fugue states, and episodic wandering.
Historical cases
1846, Albert Tirrell used sleepwalking as a defense against charges of murdering Maria Bickford, a prostitute living in a Boston brothel.
1961, Sergeant Willis Boshears confessed to strangling a local woman named Jean Constable in the early hours on New Years Day 1961, but claimed that he was asleep and only woke to realize what he had done. He pled not guilty on the basis of being asleep at the time he committed the offence and was acquitted.
In 1981, Steven Steinberg of Scottsdale, Arizona was accused of killing his wife and acquitted on the grounds of temporary insanity.
1991, R v Burgess: Burgess was accused of hitting his girlfriend on the head with a wine bottle and then a video tape recorder. He was found not guilty at Bristol Crown Court, by reason of insane automatism.
1992, R. v. Parks: Parks was accused of killing his mother-in-law and attempting to kill his father-in-law. He was acquitted by the Supreme Court of Canada.
1994, Pennsylvania v. Ricksgers: Ricksgers was accused of killing his wife. He was sentenced to life in prison without parole.
1999, Arizona v. Falater: Scott Falater, of Phoenix, Arizona, was accused of killing his wife. The court concluded that the murder was too complex to be committed while sleepwalking. Falater was convicted of first-degree murder and sentenced to life with no possibility of parole.
2001, California v. Reitz: Stephen Reitz killed his lover, Eva Weinfurtner. He told police he had no recollection of the attack but he had "flashbacks" of believing he was in a scuffle with a male intruder. His parents testified in court that he had been a sleepwalker from childhood. The court convicted Reitz of first-degree murder in 2004.
In 2001, Antonio Nieto murdered his wife and mother-in-law and attempted to murder his daughter and son, before being disarmed. Nieto claimed to have been asleep during the attack and dreaming that he was defending himself against aggressive ostriches. However, his children stated that he had recognized them and had told his son to not turn on the lights because their mother (gravely injured already) was sleeping. In 2007, Nieto was sentenced to 10 years internment in a psychiatric hospital and ordered to pay 171,100 euros as compensation to the victims.
Jules Lowe confessed to causing the death of his father Edward in 2004, but did not remember committing the act. Jules used automatism as his defense, and was found not guilty by reason of insanity and detained indefinitely in a secure hospital. He was released after ten months.
Brian Thomas was accused of killing his wife in 2008 while dreaming that he was fighting off intruders. He was freed in 2009 by a judge, who found him not guilty of murder.
| Biology and health sciences | Mental disorders | Health |
249509 | https://en.wikipedia.org/wiki/S-Bahn | S-Bahn | The S-Bahn ( , ), , is a hybrid urban–suburban rail system serving a metropolitan region predominantly in German-speaking countries. Some of the larger S-Bahn systems provide service similar to rapid transit systems, while smaller ones often resemble commuter or even regional rail systems. The name S-Bahn derives from (), (, not to be confused with the present-day Stadtbahn) or ().
Similar systems in Austria and German-speaking Switzerland are known as S-Bahn as well. In Belgium, it is known as S-Trein (Flemish) or Train S (French). In Denmark, they are known as S-tog , and in the Czech Republic as Esko or S-lines. In Milan, they are known as Linee S.
S-Bahn is also a treated as a train category in several European countries.
Characteristics
There is no complete definition of an S-Bahn system. S-Bahn are, where they exist, the most local type of passenger train service that stops at all existing stations on mainline networks inside and around a city (while other mainline trains only call at major stations).
They are slower than regional mainline trains, but usually serve as fast crosstown services within the city. The Copenhagen S-tog for example goes up to , faster than most urban heavy rail and mass transit. S-Bahn trains generally serve the hinterland of a certain city, rather than connecting different cities, although in high population density areas a few exceptions to this exist. A good example of such an exception is the Rhine-Ruhr S-Bahn, which interconnects the cities, towns and suburbs of the Ruhr, a large urban agglomeration.
Many larger S-Bahn systems have at their core a corridor of exclusive trackage that individual suburban branches feed into, creating a high frequency trunk corridor. In many cases, this central corridor is a dedicated grade-separated line in the city centre with close stop spacing and a high frequency, similar to metro systems. A good example of this is Berliner Stadtbahn in Berlin's S-Bahn, which is regarded as a tourist attraction.
Outside of the city center, most S-Bahn systems are entirely built on older local railways, or in some cases parallel to an existing dual track railway. Most use existing local mainline railway trackage, but a few branches and lines can be purpose-built S-Bahn lines. S-Bahn trains typically use overhead lines or a third rail for traction power. In Hamburg both methods are used, depending on which line is powered. In smaller S-Bahn systems and suburban sections of larger ones, trains typically share tracks with other rail traffic, with the Berlin S-Bahn, Hamburg S-Bahn and Copenhagen S-train being notable exceptions.
Further out from the central parts of a city the individual services branch off into lines where the distances between stations can exceed 5 km, similar to commuter rail. This allows the S-Bahn to serve a dual transport purpose: local transport within a city centre and suburban transport between suburbs and central boroughs of larger cities. Frequencies vary wildly between systems, with headways ranging from 2 minutes in the core sections of large networks to 30 or even 60 minutes in remote sections of the network, at off-peak times and in smaller systems.
The rolling stock typically used for S-Bahn systems reflects its hybrid purpose. The interior is designed for short journeys with provision for standing passengers, but may have more space allocated to comfortable seating than metro trains.
Integration with other local transport for ticketing, connectivity and easy interchange between lines or other systems like metros is typical for the S-Bahn. Where both S-Bahn and metro exist, the number of interchange stations between the two systems is substantial, with metro tickets being valid on S-Bahn services and vice versa. The S-Bahn Mitteldeutschland constitutes the main local railway system for Leipzig but also connects to Halle, where a few stations are located. The Rostock S-Bahn is an example of a smaller S-Bahn system.
Etymology
Germany, Austria and Switzerland
The name S-Bahn is an abbreviation of the German Stadtschnellbahn ("city rapid railway") and was introduced in December 1930 in Berlin. The name was introduced at the time of the reconstruction of the suburban commuter train tracks— the first section to be electrified was a section of the Berlin–Stettin railway from Berlin Nordbahnhof to Bernau bei Berlin station in 1924, leading to the formation of the Berlin S-Bahn.
The main line Berliner Stadtbahn ("Berlin city railway") was electrified with a 750 volt third rail in 1928 (some steam trains ran until 1929) and the circle line Berliner Ringbahn was electrified in 1929. The electrification continued on the radial suburban railway tracks along with the timetable moving to a rapid transit model with no more than a 20-minute headway per line where a number of lines overlapped on the main line. The system peaked during the 1936 Summer Olympics in Berlin with trains scheduled at least every 2 minutes.
The idea of heavy rail rapid transit was not unique to Berlin. Hamburg had an electric railway between the central station (Hauptbahnhof) and Altona which opened in 1906, and in 1934 the system adopted the S-Bahn label from Berlin. In the same year in Denmark, Copenhagen's S-tog opened its first line. In Austria, Vienna had its Stadtbahn main line electrified in 1908 and also introduced the term Schnellbahn ("rapid railway") in 1954 for its planned commuter railway network, which started operations in 1962. The S-Bahn label was sometimes used as well, but the name was only switched to S-Bahn Wien in 2005.
As for Munich, in 1938 the Nazi government broke ground for an S-Bahn-like rapid transport system in Lindwurmstraße near what is now Goetheplatz station on line U6. The system was supposed to run through tunnels in the city centre. The planning process mainly consisted of the bundling and interconnecting of existing suburban and local railways, plus the construction of a few new lines. Plans and construction work - including the building shell of Goetheplatz station - came to a very early halt during World War II and were not pursued in its aftermath. Very extensive nowadays, Munich's existing S-Bahn system, together with the first two U-Bahn lines, only began to operate prior to the 1972 Summer Olympics.
The term S-Bahn was a registered wordmark of Deutsche Bahn until 14 March 2012, when, at the request of a transportation association, the Federal Patent Court of Germany ordered its removal from the records of the German Patent and Trade Mark Office. Prior to this Deutsche Bahn collected a royalty of 0.4 cents per train kilometer for the use of the term.
Denmark
The "S" stood for "station". Just before the opening of the first line in the Copenhagen S-train network, the newspaper Politiken on 17 February 1934 held a competition about the name, which in Danish became known as Den elektriske enquete or "The electrical survey" (as the Copenhagen S-trains would become the first electrical railways in Denmark). But since an "S" already was put up at all the stations, weeks before the survey, the result became S-tog which means "S-train". This was also just a few years after the S-trains had opened in Berlin and Hamburg. Today the Copenhagen S-trains uses six lines and serves 86 stations, 32 of them are located inside the (quite tiny) municipality borders. Each line uses 6 t.p.h (trains per hour) in each direction, with exception of the (yellow) F-line. The F-line has departures in each direction every five minutes, or 12 t.p.h. service .
History
The history of the S-Bahn in present-day Germany begins in the 19th century in Prussia.
Early steam services
In 1882, the growing number of steam-powered trains around Berlin prompted the Prussian State Railway to construct separate rail tracks for suburban traffic. The Berliner Stadtbahn connected Berlin's eight intercity rail stations which were spread throughout the city (all but the Stettiner Bahnhof which today is a pure S-Bahn station known as Berlin Nordbahnhof; as the city Stettin today is Polish city Szczecin). A lower rate for the newly founded Berliner Stadt-, Ring- und Vorortbahn (Berlin City, Circular and Suburban Rail) was introduced on 1 October 1891. This rate and the growing succession of trains made the short-distance service stand out from other railways.
The second suburban railway was the Hamburg-Altonaer Stadt- und Vorortbahn connecting Hamburg with Altona and Blankenese. The Altona office of the Prussian State Railway established the electric powered railway in 1906.
Electricity
The beginning of the 20th century saw the first electric trains, which in Germany operated at 15,000 V on overhead lines. The Berliner Stadt-, Ring- und Vorortbahn instead implemented direct current multiple units running on 750 V from a third rail. In 1924, the first electrified route went into service. The third rail was chosen because it made both the modifications of the rail tracks (especially in tunnels and under bridges) and the side-by-side use of electric and steam trains easier.
To set it apart from the subterranean U-Bahn, the term S-Bahn replaced Stadt-, Ring- und Vorortbahn in 1930.
The Hamburg service had established an alternating current line in 1907 with the use of multiple units with slam doors. In 1940 a new system with 1200 V DC third rail and modern electric multiple units with sliding doors was integrated on this line (on the same tracks). The old system with overhead wire remained up to 1955. The other lines of the network still used steam and later Diesel power.
In 1934, the Hamburg-Altonaer Stadt- und Vorortbahn was renamed as S-Bahn.
Systems by country
Austria
The oldest and largest S-Bahn system in Austria is the Vienna S-Bahn, which predominantly uses non exclusive rails tracks outside of Vienna. It was established in 1962, although it was usually referred to as Schnellbahn until 2005. The white "S" on a blue circle used as the logo is said to reflect the layout of the central railway line in Vienna. However, it has now been changed for a more stylized version that is used all through Austria, except Salzburg. The rolling stock was blue for a long time, reflecting the logo colour, but red is used uniformly for nearly all local traffic today.
In 2004, the Salzburg S-Bahn went into service as the first Euroregion S-Bahn, crossing the border to the neighbouring towns of Freilassing and Berchtesgaden in Bavaria. The network is served by three corporations: the Berchtesgadener Land Bahn (BLB)(S4), the Austrian Federal Railways (German: Österreichischen Bundesbahn / ÖBB)(S2 and S3) and the Salzburger Lokalbahn (SLB)(S1 and S11) and . The Salzburg S-Bahn logo is only different one, it is a white S on a light blue circle.
In 2006 the regional train line in the Rhine Valley in the state of Vorarlberg has been renamed to Vorarlberg S-Bahn. It is a three lines network, operated by the Montafonerbahn and the ÖBB. It was later expanded. Presently, a frequent service, the S1, operates between to (D) via . In addition, an hourly service, S3 (ÖBB), connects Bregenz with St. Margrethen (CH), and another service (S2) operates between Feldkirch (A), Schaan (FL) and Buchs SG (CH). The Montafonerbahn runs the S4.
The S-Bahn Steiermark has been inaugurated in December 2007 in Styria, built to connect its capital city Graz with the rest of the metropolitan area, currently the following lines are active: S1, S11, S3, S31, S5, S51, S6, S61, S7, S8 and S9. The network is operated by three railway companies: the Graz-Köflacher Bahn (GKB) (lines: S6, S61 and S7), the ÖBB (lines: S1, S3, S5, S51, S8 and S9) and the Steiermärkische Landesbahnen (StB) (lines: S11 and S31).
In December 2007 as well the Tyrol S-Bahn opened, running from Hall in Tirol in the east to Innsbruck Central Station and Telfs in the west and from Innsbruck to Steinach am Brenner. Class 4024 EMUs are used as rolling stock on this network.
In 2010 the S-Bahn Kärnten was opened in the state of Carinthia and currently consists of 4 lines operated by ÖBB.
The youngest network is the S-Bahn Oberösterreich in the Greater Linz area of the state of Upper Austria, which was inaugurated in December 2016. It is a 5 line system operated by Stern und Hafferl and the ÖBB.
Belgium
Since 2015, the trains of the Brussels Regional Express Network (French: Réseau Express Régional Bruxellois, RER; Dutch: Gewestelijk Expresnet, GEN) of the NMBS/SNCB belong to train category S and are referred to as S train (Dutch: 'S-trein', French: train S, German: 'S-Züge'). In 2018, local trains of NMBS/SNCB in and around Antwerp, Ghent, Liège and Charleroi changed to the train category S train as well.
Czech Republic
In the Czech Republic, integrated commuter rail systems exist in Prague and Moravian-Silesian Region. Both systems are called Esko, which is how S letter is usually called in Czech. Esko Prague has been operating since 9 December 2007 as a part of the Prague Integrated Transport system. Esko Moravian-Silesian Region began operating on 14 December 2008 as a part of the ODIS Integrated Transport system serving the Moravian-Silesian Region. Both systems are primarily operated by České dráhy. Several shorter lines are operated by other companies.
Denmark
Copenhagen S-train connects the city centre, other inner and outer boroughs and suburbs with each other. The average distance between stations is 2.0 km, shorter in the city core and inner boroughs, longer at the end of lines that serve suburbs. Of the 86 stations, 32 are located within the central parts of the city. Some stations are located around 40 km from Copenhagen city centre. For this reason the fares vary depending on distances. The one-day passes which the tourists buy are valid only in the most central parts of the S-train system. On weekdays each line has a departure every 10 minutes with the exception of the F-line, on which a train departs every five minutes. Where several lines converge on a common piece of track there could be as many as 30 trains per hour in each direction. On Sundays the seven lines are reduced to four lines, but all stations are served at least every 10 minutes. The three railway stations at Amager have a local service that is the equivalent of the S-trains.
The Copenhagen Metro opened in 2002 as a complement to the already existing S-train system. Copenhagen's S-train system is the only one in the country. Outside Denmark, in cities where both exist, is it far from unusual that a metro system later has been complemented with S-trains. The branch towards Køge (the southernmost S-train station in Copenhagen's S-network) has a rather unique history, as it was built in the 1970s where no previous railway ever had existed.
France
Although not called a "Train S" in French, the Paris RER has the characteristic structure of an S-Bahn system, with branches going to distant suburbs and sharing a common corridor in the city center. It is also called "S-Bahn" in some German-language signs.
Germany
The trains of the Berlin and Hamburg S-Bahn systems ran on separate tracks from the beginning. When other cities started implementing their systems in the 1960s, they mostly had to use the existing intercity rail tracks, and they still more or less use such tracks.
The central intercity stations of Frankfurt, Leipzig, Munich and Stuttgart are terminal stations, so all four cities have monocentric S-Bahn networks. The S-Bahn trains use as their core segment a tunnel under the central station and the city centre (e.g. Munich S-Bahn Stammstrecke and the upcoming Zweite Stammstrecke).
The high number of large cities in the Ruhr area promotes a polycentric network connecting all cities and suburbs. The S-Bahn Rhein-Ruhr, as it is called, features few tunnels, and its routes are longer than those of other networks. The Ruhr S-Bahn is the only S-Bahn network to be run by more than one corporation in Germany, and the Salzburg S-Bahn holds a similar distinction in Austria. Most Swiss S-Bahn systems are multi-corporation networks, however.
Most German S-Bahn networks have a unique ticket system, separated from the Deutsche Bahn rates, instead connected to the city ticket system used for U-bahns and local buses. The S-Bahn of Hanover, however, operates under five different rates due to its large expanse.
One S-Bahn system is no longer in operation: the Erfurt S-Bahn which operated from 1976 until 1995 and was an single-line diesel-powered system which consisted of four stations from Erfurt Central Station to Erfurt Berliner Straße station in the then newly built northern suburbs of Erfurt.
There are several S-Bahn or S-Bahn-like systems in planning, such as the Augsburg S-Bahn (network plan), the Lübeck S-Bahn (network plan) and the tri-country Bodensee S-Bahn.
The Stadtbahn Karlsruhe (a tram-train network) uses the green "S" logo for stations in the outskirts and has its lines indicated by an "S" in front of the line number, but does not refer to itself as S-Bahn. The logo also can't be found on the trains, contrary to most other systems where it's placed somewhere on the sides or at the front of the trains. A new city-centre tunnel opened at the end of 2021, however the blue U-Bahn logo is not used either for it. To mark those tunnel stations, a yellow U is used, which is unique and can only be found there.
Despite their names, the Ortenau S-Bahn (Offenburg) and the Danube-Iller Regional S-Bahn (Ulm/Neu-Ulm, opened 2020) are Regionalbahn services.
The following networks are currently in operation:
Liechtenstein
The only railway line passing through the Principality of Liechtenstein is the Feldkirch–Buchs railway line, which connects with the Austrian rail network in and with the Swiss network in . In June 2008, the Swiss canton of St. Gallen, the Austrian state of Vorarlberg, and the Principality of Liechtenstein signed an agreement for a project to upgrade this line (and the surrounding ones) and to increase the rail traffic. The project, named was approved by Liechtenstein and Austria in a Letter of Intent signed in April 2020 and under that plan, it was to be fully realised by 2027 and would have cost an estimated €187 million. That plan was however rejected by 62.3% of Liechtenstein voters in a referendum on 30 August 2020.
As of the December 2023 timetable change, an S-Bahn service, the S2 of Vorarlberg S-Bahn, operates between Feldkirch (A), Schaan (FL) and Buchs SG (CH). There are three operational railway stations in Liechtenstein along the Feldkirch–Buchs line: (which serves the capital Vaduz), and . A fourth station, , was closed in 2013.
Switzerland
S-Bahn is also used in the German-speaking part of Switzerland. Swiss French networks use the term RER with line numbers prefixed with an R, e.g. as R2, except for the Léman Express in Greater Geneva that uses the prefix L followed by the line number ("L" for "Léman-Express"), e.g. L2. S-Bahn-style services in the Italian and Romansh speaking parts of Switzerland also use, like the Milan suburban system, the "S" prefix, although in Italian such networks are called rete celere () instead of S-Bahn.
The oldest network in Switzerland is the Bern S-Bahn, which was established in stages from 1974 onward and has adopted the term S-Bahn since 1995. It is also the only one in Switzerland to use a coloured "S" logo. In 1990, the Zürich S-Bahn, went into service. As of 2022, this network comprises 32 services, covering a large area in Switzerland (and parts of southern Germany). Further S-Bahn services were set up in the course of the Bahn 2000 initiative in Central Switzerland (a collaborative network of S-Bahn Luzern and Stadtbahn Zug), and Eastern Switzerland (S-Bahn St. Gallen).
The Basel trinational S-Bahn services the Basel metropolitan area, thus providing cross-border transportation into both France and Germany. A tunnel connecting Basel's two large intercity stations (Basel Badischer Bahnhof and Basel SBB) is planned as Herzstück Regio-S-Bahn Basel (lit. heart-piece Regio-S-Bahn Basel).
An international S-Bahn network also existsts across the Swiss-Italian border, in the Swiss Canton of Ticino and the Italian state of Lombardy. Services are operated by Treni Regionali Ticino Lombardia (TILO), a joint venture between Italian railway company Trenord and Swiss Federal Railways (SBB CFF FFS).
The RER Vaud of Lausanne and the Léman Express of Geneva serve the area around Lake Geneva (fr. Lac Léman). The Léman express network expands across the Swiss-French border. It is the largest cross-country S-Bahn network of Europe. Léman express was launched in December 2019 and is operated by Swiss Federal Railways (SBB CFF FFS) and SNCF.
Another transborder network for the Lake Constance (Bodensee) area, connecting up to four nations, is under discussion. This network would extend across the German states Baden-Württemberg and Bavaria, the Austrian state Vorarlberg, the Principality of Liechtenstein (S-Bahn FL.A.CH), and the Swiss cantons of Appenzell Ausserrhoden, Appenzell Innerrhoden, Schaffhausen, St. Gallen and Thurgau. Possible names are Bodensee-S-Bahn and Alpenrhein-Bahn. Presently, the Bodensee S-Bahn only operates services around Lake Constance in Austria, Germany and Switzerland (without Liechtenstein). It includes, among others, the S14 and S44 services of St. Gallen S-Bahn, which both connect Konstanz (D) with Kreuzlingen and Weinfelden (CH). Since 2022, some S7 services continue from Rorschach (CH) to Bregenz (A) and Lindau-Reutin (D). Additional transborder services are planned.
The Chur S-Bahn provides services around Chur, the capital of the alpine Canton of Graubünden (Grisons) in south-eastern Switzerland.
The Aargau S-Bahn is a small network that services stations in the cantons of Aargau, Lucerne and Bern.
The RER Fribourg is an S-Bahn-style service centered at Fribourg/Freiburg and Bulle in the Canton of Fribourg, and extending into the cantons of Neuchâtel and Vaud.
Two unnumbered S-Bahn services (designated only with an "S"), one between Schaffhausen and Erzingen (D), running on railway tracks owned by Deutsche Bahn (DB), and one between Schaffhausen and Jestetten (D), opened in 2013. They are operated by SBB GmbH and THURBO, respectively. Since December 2022, the Schaffhausen–Singen am Hohentwiel line is also serviced by SBB GmbH As of the December 2023 timetable change, the three services of Schaffhausen S-Bahn are numbered S62, S64 and S65.
Additionally, there are services designated "S" that are not part of any formal S-Bahn network. These include the S20, S21, and S22 operated by Swiss Federal Railways in Solothurn or the S27 operated by Südostbahn (SOB) between Siebnen-Wangen and Ziegelbrücke.
Swiss S-Bahn services are operated mostly by the Swiss Federal Railways (SBB CFF FFS) but also by private railway companies, such as Appenzeller Bahnen (AB), BLS AG, Forchbahn (FB), Regionalverkehr Bern-Solothurn (RBS), Rhätische Bahn (RhB), Sihltal Zürich Uetliberg Bahn (SZU), Südostbahn (SOB) or Zentralbahn (ZB).
Rail transport in Switzerland, including S-Bahn systems, is noteworthy for its coordination between services due to the clock-face schedule. Due to the proximity of the various S-Bahn systems in Switzerland, services of one network often offer connections to services of neighboring networks. S-Bahn services are used by commuters and tourists (some services call nearby tourist attractions, such as the Rhine Falls or the Swiss Museum of Transport).
| Technology | Germany | null |
249530 | https://en.wikipedia.org/wiki/Plains%20zebra | Plains zebra | The plains zebra (Equus quagga, formerly Equus burchellii) is the most common and geographically widespread species of zebra. Its range is fragmented, but spans much of southern and eastern Africa south of the Sahara. Six or seven subspecies have been recognised, including the quagga which was thought to be a separate species. More recent research supports variations in zebra populations being clines rather than subspecies.
Plains zebras are intermediate in size between the larger Grévy's zebra and the smaller mountain zebra and tend to have broader stripes than both. Great variation in coat patterns exists between clines and individuals. The plains zebra's habitat is generally, but not exclusively, treeless grasslands and savanna woodlands, both tropical and temperate. They generally avoid desert, dense rainforest and permanent wetlands. Zebras are preyed upon by lions and spotted hyenas, Nile crocodiles and, to a lesser extent, leopards, cheetahs and African wild dogs.
Plains zebras are a highly social species, forming harems with a single stallion, several mares and their recent offspring; bachelor groups also form. Groups may come together to form herds. The animals keep watch for predators; they bark or snort when they see a predator and the harem stallion attacks predators to defend his harem.
The plains zebra remains common in game reserves, but is threatened by human activities, such as hunting for its meat and hide, as well as competition with livestock and encroachment by farming on much of its habitat. The loss of open grasslands due to woody plant encroachment increases predation risk and therewith habitat. Plains zebra are listed as near threatened by the IUCN as of 2016. The species population is stable and not endangered, though populations in most countries have declined sharply.
Taxonomy
The plains zebra was formally classified by British zoologist John Edward Gray in 1824 as Equus burchellii. After the quagga, described by Pieter Boddaert in 1785, was found to be the same species in the 21st century, the plains zebra was reclassified as Equus quagga due to the principle of priority. The plains zebra and mountain zebra were traditionally placed in the subgenus Hippotigris, in contrast to Grévy's zebra, which was considered the sole species of the subgenus Dolichohippus; however, recent (2013) phylogenetic evidence finds that plains zebras are more closely related to Grévy's zebras than mountain zebras. Groves and Bell (2004) place all three species in the subgenus Hippotigris, and zebras appear to be a monophyletic lineage. In areas where plains zebras are sympatric with Grévy's zebras, finding them in the same herds is not unusual, and fertile hybrids occur.
Subspecies
In their 2004 study of cranial and pelage differences between specimens, Groves and Bell found support for the division of the plains zebra into six subspecies:
*Sometimes a seventh subspecies is recognised.
Burchell's zebra was thought to have been hunted to extinction. However, Groves and Bell concluded that "the extinct true Burchell's zebra is a phantom". Careful study of the original zebra populations in Zululand and Eswatini and of skins harvested on game farms in Zululand and Natal has revealed that a certain small proportion shows similarity to what now is regarded as typical burchellii. The type localities of the subspecies Equus quagga burchellii and Equus quagga antiquorum are so close to each other that the two are in fact one and that, therefore, the older of the two names should take precedence over the younger. They suggested that the correct name for the subspecies must be burchellii, not antiquorum.
A 2005 genetic study confirmed the quagga being the same species as the plains zebra. It showed that the quagga had little genetic diversity and that it diverged from the other plains zebra subspecies only 120,000–290,000 years ago, during the Pleistocene and, possibly, the penultimate glacial maximum. Its distinct coat pattern may have evolved rapidly because of geographical isolation and/or adaptation to a drier environment. In addition, plains zebra subspecies tend to have less striping the further south they live, and the quagga was the most southern-living of them all.
The simplified cladogram below is based on the 2005 analysis (some taxa shared haplotypes and could therefore not be differentiated):
A 2018 DNA study found no evidence for a subspecies structure in plains zebras but, instead, observed a north–south genetic continuum. Modern plains zebra populations appear to have originated from Southern Africa around 370,000 years ago with plains zebras in Uganda, the most northern population, being the most distinct.
Physical description
The plains zebra stands at a height of with a head-body length of and a tail length of . Males weigh while females weigh . The species is intermediate in size between the larger Grévy's zebra and the smaller mountain zebra. It is dumpy bodied with relatively short legs and a skull with a convex forehead and a somewhat concave nose profile. The neck is thicker in males than in females. The ears are upright and have rounded tips. They are shorter than in the mountain zebra and narrower than in Grévy's zebra. As with all wild equids, the plains zebra has an erect mane along the neck and a tuft of hair at the end of the tail. The body hair of a zebra is , shorter than in other African ungulates.
Like all zebras, they are boldly striped in black and white and no two individuals look exactly alike. Compared to other species, the plains zebra has broader stripes. The stripes are vertical on the fore part of the body, and tend towards the horizontal on the hindquarters. Northern zebra populations have narrower and more defined striping; southern populations have varied but lesser amounts of striping on the under parts, the legs and the hindquarters. Southern populations also have brown "shadow" stripes between the black and white colouring. These are absent or poorly expressed in northern zebras. The natal coat of a foal is brown and white and the brown darkens with age.
Various abnormalities of the patterns have been documented in plains zebras. Melanistic zebras have high concentrations of dark stripes on the torso but low concentrations on the legs. "Spotted" individuals display interruptions in black striping patterns. There have even been morphs with white spots on dark backgrounds. Striping abnormalities have been linked to inbreeding. Albino zebras have been recorded in the forests of Mount Kenya, with the dark stripes being blonde. The quagga had brown and white stripes on the head and neck, brown upper parts and a white belly, tail and legs.
Ecology
Range and habitat
The plains zebra's range stops short of the Sahara from South Sudan and southern Ethiopia extending south along eastern Africa, as far as Zambia, Mozambique, and Malawi, before spreading into most southern African countries. They may have lived in Algeria in the Neolithic era. Plains zebras generally live in treeless grasslands and savanna woodlands, but can be found in a variety of habitats, both tropical and temperate. However, they are generally absent from deserts, dense rainforests, and permanent wetlands. They generally prefer Acacieae woodlands over Commiphora. They are water-dependent and live in more mesic environments than other African equids. They seldom wander from a water source. Zebras also live in elevations from sea level to on Mount Kenya.
Depending on the population, zebra herds may be sedentary, being highly dense with small ranges, or migratory, being less populated with separate, extensive dry and wet home ranges. When migrating, zebras appear to rely on some memory of the locations where foraging conditions were best and may predict conditions months before their arrival.
The loss of open grasslands due to woody plant encroachment increases predation risk and therewith habitat.
Diet and predation
Plains zebras primarily feed on grass; preferred species being Themeda triandra, Cynodon dactylon, Eragrostis superba and Cenchrus ciliaris. Zebra sometimes browse or dig for corms and rhizomes during the dry season. They appear to partial to eating scorched Colophospermum mopane and Pterocarpus rotundifolius, consuming both the leaves and twigs.
Plains zebras are adapted for grazing on both long, tough grass stems and newly emerging short grass. In some areas, it rarely feeds below to ground level. It ranges more widely than many other species, even into woodlands, and it is often the first grazing species to appear in a well-vegetated area. The flexible upper lip allows them to push plant material between the incisors to cut. Zebras have a less efficient digestive system than ruminants but food passage is twice as fast. Thus, zebras are less selective in foraging, but they do spend much time eating. The zebra is a pioneer grazer and prepares the way for more specialised grazers such as blue wildebeests and Thomson's gazelles.
The plains zebra's major predators are lions and spotted hyenas. Lions are most successful when targeting lone individuals, usually an old male while hyenas chase and isolate an individual from the group, usually a female or foal. Nile crocodiles also prey on zebras when they are near water. Less common predators include leopards, cheetahs and African wild dogs, which mostly hunt foals. When in the presence of a lion, zebras remain alert and stand in a semi-circle at as much as and no less than . Stallions sometimes try to drive lions away with bluff charges. By contrast, zebras may approach cheetahs and wild dogs and a single hyena is allowed to come within a few metres. To escape from predators, an adult zebra can run at . When being hunted by hyenas or wild dogs, a zebra harem stays close together and cooperates to protect threatened members, particularly the young. The harem stallion goes on the offensive and attacks the dogs or hyenas.
Behaviour
Plains zebras are nomadic and non-territorial, home ranges vary from to , depending on the area and if the population is migratory. They are more active during the day and spend most of their time feeding. Other activities include dust bathing, rubbing, drinking and intermittent resting which is very brief. At night, zebra activity is subdued except when threatened by predators. They may rest or sleep laying down, while one individual keeps guard.
Social structure
The plains zebra is highly social and usually forms small family groups called harems, which consist of a single stallion, several mares and their recent offspring. The adult membership of a harem is highly stable, typically remaining together for months to years. Groups of all-male "bachelors" also exist. These are stable groups of up to 15 males with an age-based hierarchy, led by a young male. These males stay in their groups until they are ready to start a harem. The bachelors prepare for their adult roles with play fights and greeting/challenge rituals, which take up most of their activities.
Multiple harems and bachelor groups come together to form larger herds of hundreds of animals, especially during migrations. Plains zebras are unusual among harem-holding species in forming these groups. In addition, pairs of harems may create temporarily stable subgroups within a herd, allowing individuals to interact with those outside their group. Among harem-holding species, this has only been observed in primates such as the gelada and the hamadryas baboon. Bachelor groups tend to be at the periphery of herds and when the herd moves, the bachelors trail behind.
Stallions form and expand their harems by abducting young mares from their natal harems. When a mare reaches sexual maturity and has her first oestrous cycle, she attracts the attention of nearby stallions, both bachelors and harem leaders. Her family stallion (likely her father) chases off or fights stallions attempting to abduct her. Even after a young mare is isolated from her natal harem, the fight over her continues until her oestrous cycle is over and it starts again with the next oestrous cycle. It is rare that the mare's original abductor keeps her for long. When the mare finally ovulates, the male that impregnates her keeps her for good. Thus, the mare becomes a permanent member of a new harem. Oestrus in a female becomes less noticeable to outside males as she gets older, hence competition for older females is virtually nonexistent.
Mares exist in a hierarchy, with the alpha female being the first to mate with the harem stallion and being the one to lead the group. When new mares are added to the group, they are met with hostility by the other mares. Thus, the harem stallion must shield the new mares until the aggression subsides. The most recently added females rank lowest. Females that become unfit or weak may drop in their rank, though. The female membership of a harem stays intact, even if a new stallion takes over. During herd gatherings, family stallions may be cordial towards each other, while the mares are less tolerant.
A stallion defends his harem from other males. When challenged, the stallion issues a warning to the invader by rubbing nose or shoulder with him. If the warning is not heeded, a fight breaks out. Zebra fights often become very violent, with the animals biting at each other's necks, heads, or legs, wrestling to the ground and occasional kicking. Sometimes, a stallion lies still on the ground as if surrendering, but once the other male lets up, he strikes and continues the fight. Most fighting occurs over young mares in oestrus and as long as a harem stallion is healthy, he usually is not challenged. Only unhealthy stallions have their harems taken over and even then, the new stallion gradually takes over, pushing the old one out without a fight.
Communication
At least six different calls have been documented for the plains zebra. One is its distinctive, high-pitched, contact call (commonly called "barking") heard as "a-ha, a-ha, a-ha" or "kwa-ha, kaw-ha, ha, ha" also transcribed as "kwahaah", or "oug-ga". The species name quagga is derived from the Khoikhoi word for "zebra" and is onomatopoeic for its call. When a predator is sighted, a zebra makes a two-syllable alarm call. A loud snort is made when moving in cover of potential danger. When in contentment, a zebra makes a more drawn-out snort. Males make a short, high-pitched squeal when hurt, and foals emit a drawn-out wail when in distress.
Two main facial expressions are made by zebras; the greeting and threat. In both cases, the lips are pulled back and chewing motions are made. Greeting involves the ears sticking up and directing forward; while the threat involves the ears down. Zebras strengthen their social bonds with grooming. Members of a harem nip and scrape along the neck, withers, and back with their teeth and lips. Mothers and foals groom the most often, followed by siblings. Grooming shows social status and eases aggressive behaviour.
Reproduction and parenting
The stallion mates with all his mares. Males exhibit the flehmen response to test for female receptivity, which involves the upper lip curling back to smell for urine (via the vomeronasal organ). The female signals her readiness for copulation by straddling her legs and raising her tail. The gestation period lasts around a year, and a single young is produced. Mares may give birth to one foal every twelve months. The birthing peak is during the rainy season. A mare gives birth within the vicinity of her group and while lying down on her side. The newborn foal weighs and the afterbirth is rarely consumed.
A newborn is capable of standing almost immediately and starts to eat grass within a week. Early on, a mother zebra keeps any other zebra away from her foal, including the stallion, the other mares, and even her previous offspring. Later, though, they all bond. Within the group, a foal has the same rank as its mother. The stallion is generally intolerant of foals that are not his, and zebras may practice infanticide and feticide.
Mortality for foals is high in their first year of life and is usually caused by predation. However, zebra young are afforded more protection than those of species like wildebeest and hartebeest. A foal is usually weaned at around eleven months, but may suckle for longer. Females reach puberty before three years, and males after five or six. Young male zebras eventually leave their family groups as the relationship with their mothers fades after the birth of a sibling. The young stallion then seeks out other young stallions for company. Young females may stay in the harem until they are abducted by another stallion. Plains zebras have an average lifespan of 25 years in the wild.
Human interactions
Conservation
In 2016, the plains zebra was classified as near threatened by the IUCN. As of that year, the total population is estimated to be around 500,000 individuals. The species remains common throughout its range but has experienced population declines in 10 of the 17 countries where it is native. They are stable in Ethiopia, Malawi, and South Africa and possibly Angola; stable or increasing in Mozambique, Namibia and Eswatini; and decreasing in Botswana, DR Congo, Kenya, Rwanda, Somalia, South Sudan, Tanzania, Uganda, Zambia, and Zimbabwe. They are extinct in Burundi, Lesotho and possibly Somalia.
Zebras are threatened by hunting for their hide and meat, and habitat change from farming. They also compete with livestock for food, and fencing blocks migration routes. Civil wars in some countries have also caused declines in zebra populations. The zebra can be found in numerous protected areas across its range, including the Serengeti National Park in Tanzania, Tsavo and Masai Mara in Kenya, Hwange National Park in Zimbabwe, Etosha National Park in Namibia, and Kruger National Park in South Africa. Some stable populations live in unprotected areas.
The quagga was hunted by early Dutch settlers and later by Afrikaners to provide meat or for their skins. The skins were traded or used locally. The quagga was probably vulnerable to extinction due to its limited distribution, and it may have competed with domestic livestock for forage. The last known wild quagga died in 1878. The last captive quagga, a female in Amsterdam's Natura Artis Magistra zoo, lived there from 9 May 1867 until it died on 12 August 1883, but its origin and cause of death are unclear. In 1984, the quagga was the first extinct animal to have its DNA analysed, and the Quagga Project is trying to recreate the phenotype of hair coat pattern and related characteristics by selectively breeding Burchell's zebras.
In popular culture
Zebras have been featured in African art and culture for millennia. They have been depicted in rock art in Southern Africa (modern Botswana, Namibia and South Africa) dating from 20,000–28,000 years ago, though not as commonly as antelope species like eland. How the zebra got its stripes has been the subject of folk tales, some of which involve it being scorched by fire. The San people associated zebra stripes with water, rain and lighting due to its dazzling pattern.
The plains zebra is the national animal of the Republic of Botswana and its stripes are depicted on the country's flag. The flag stripes also represent racial harmony in the country. The zebra has also been associated with beauty and the women of various societies would paint much of their bodies in stripes. For the Shona people of Zimbabwe, the zebra is a totem animal, along with the eland, buffalo, lion and monkey. The zebra is praised in a poem as an "iridescent and glittering creature". Its stripes have symbolised the joining of male and female and at Great Zimbabwe, zebra stripes decorate what is believed to be a domba, a premarital school meant to initiate women into adulthood. In the Shona language, the name "madhuve" means "woman/women of the zebra totem" and is a given name for girls in Zimbabwe.
Zebras have also been represented in Western culture. They have been thought of as a more exotic alternative to horse; the comic book character Sheena, Queen of the Jungle is depicted riding a zebra. The film Racing Stripes features a captive zebra ostracised from the horses and end up being ridden by a rebellious girl. In the film Fantasia, two centaurs are depicted being half human and half zebra, instead of the typical half human and half horse. Zebras have been featured as characters in other animated films like Khumba, The Lion King and the Madagascar films.
| Biology and health sciences | Equidae | Animals |
249613 | https://en.wikipedia.org/wiki/Continental%20climate | Continental climate | Continental climates often have a significant annual variation in temperature (warm to hot summers and cold winters). They tend to occur in central and eastern parts of the three northern-tier continents (North America, Europe, and Asia), typically in the middle latitudes (40 to 55 or 60 degrees north), often within large landmasses, where prevailing winds blow overland bringing some precipitation, and temperatures are not moderated by oceans. Continental climates occur mostly in the Northern Hemisphere due to the large landmasses found there. Most of northeastern China, eastern and southeastern Europe, much of Russia south of the Arctic Circle, central and southeastern Canada, and the central and northeastern United States have this type of climate. Continentality is a measure of the degree to which a region experiences this type of climate.
In continental climates, precipitation tends to be moderate in amount, concentrated mostly in the warmer months. Only a few areas—in the mountains of the Pacific Northwest of North America and in Iran, northern Iraq, adjacent Turkey, Afghanistan, Pakistan, and Central Asia—show a winter maximum in precipitation. A portion of the annual precipitation falls as snowfall, and snow often remains on the ground for more than a month. Summers in continental climates can feature thunderstorms and frequent hot temperatures; however, summer weather is somewhat more stable than winter weather. Continental climates are considered as temperate climate varieties due to their location in the temperate zones, but are classified separately from other temperate climates in the Köppen climate classification system where they are identified by their first letter, a capital D. In the Trewartha climate classification, they are identified as Dc.
Köppen climate classification
Continental climate has at least one month averaging below and at least one month averaging above .
Dfa = Hot-summer humid continental climate; coldest month averaging below (or ), at least one month's average temperature above , and at least four months averaging above . No significant precipitation difference between seasons (neither the abovementioned set of conditions fulfilled).
Dfb = Warm-summer humid continental climate; coldest month averaging below (or ), all months with average temperatures below , and at least four months averaging above . No significant precipitation difference between seasons (neither the abovementioned set of conditions fulfilled).
Dfc = Subarctic climate; coldest month averaging below (or ) and one–three months averaging above . No significant precipitation difference between seasons (neither the abovementioned set of conditions fulfilled).
Dfd = Extremely cold subarctic climate; coldest month averaging below and one–three months averaging above . No significant precipitation difference between seasons (neither the abovementioned set of conditions fulfilled).
Dwa = Monsoon-influenced hot-summer humid continental climate; coldest month averaging below (or ), at least one month's average temperature above , and at least four months averaging above . At least ten times as much rain in the wettest month of summer as in the driest month of winter.
Dwb = Monsoon-influenced warm-summer humid continental climate; coldest month averaging below (or ), all months with average temperatures below , and at least four months averaging above . At least ten times as much rain in the wettest month of summer as in the driest month of winter.
Dwc = Monsoon-influenced subarctic climate; coldest month averaging below (or ) and one–three months averaging above . At least ten times as much rain in the wettest month of summer as in the driest month of winter.
Dwd = Monsoon-influenced extremely cold subarctic climate; coldest month averaging below and one–three months averaging above . At least ten times as much rain in the wettest month of summer as in the driest month of winter.
Dsa = Mediterranean-influenced hot-summer humid continental climate; coldest month averaging below (or ), average temperature of the warmest month above and at least four months averaging above . At least three times as much precipitation in the wettest month of winter as in the driest month of summer, and the driest month of summer receives less than .
Dsb = Mediterranean-influenced warm-summer humid continental climate; coldest month averaging below (or ), average temperature of the warmest month below and at least four months averaging above . At least three times as much precipitation in the wettest month of winter as in the driest month of summer, and the driest month of summer receives less than .
Dsc = Mediterranean-influenced subarctic climate; coldest month averaging below (or ) and one–three months averaging above . At least three times as much precipitation in the wettest month of winter as in the driest month of summer, and the driest month of summer receives less than .
Dsd = Mediterranean-influenced extremely cold subarctic climate; coldest month averaging below and one–three months averaging above . At least three times as much precipitation in the wettest month of winter as in the driest month of summer, and the driest month of summer receives less than .
Seasons
Annual precipitation in this zone is usually between and , The timing of intermediate spring-like or autumn-like temperatures in this zone vary depending on latitude and/or elevation. For example, spring may arrive as soon as March (in Northern hemisphere, September in Southern hemisphere) in the southern (in Northern hemisphere, northern in Southern hemisphere), parts of this zone or as late as May (November) in the north (south). Summers are warm or hot while winters are below freezing and sustain lots of frost.
Climatology
Continental climates exist where cold air masses infiltrate during the winter from shorter days and warm air masses form in summer under conditions of high sun and longer days. Places with continental climates are as a rule either far from any moderating effect of oceans or are so situated that prevailing winds tend to head offshore. Such regions get quite warm in the summer, achieving temperatures characteristic of tropical climates but are colder than any other climates of similar latitude in the winter.
Neighboring climates
In the Köppen climate system, these climates grade off toward temperate climates equator-ward where winters are less severe and semi-arid climates or arid climates where precipitation becomes inadequate for tall-grass prairies and shrublands. In Europe these climates may grade off into oceanic climates (Cfb) or subpolar oceanic climates (Cfc) in which the influence of cool oceanic air masses is more marked toward the west. In western and eastern Asia, and the central United States these climates grade off toward humid subtropical climates (Cfa/Cwa), subtropical highland climates (Cwb), or Mediterranean climates (Csa/Csb) to the south.
List of locations with a continental climate
The climate is continental if the 0 °C coldest-month isotherm is used, but it is temperate if the -3 °C isotherm is used.
Africa
Morocco
Imilchil (bordering Cfb)
Eurasia
Asia
Afghanistan: Chaghcharan (Dsb), Fayzabad (Dsa bordering Csa), Ghazni (Dsa bordering BSk), Maidan Shar (Dsb), Puli Alam (Dsa)
Armenia: Gyumri, Jermuk (Dsb bordering Dfb), Kapan, Sisian, Vanadzor, Yerevan (bordering BSk)
Azerbaijan: Qabala (bordering Cfa), Shamakhi (bordering Cfa)
China: Anshan, Beijing, Biru County, Changchun, Chengde, Dalian, Daqing, Hailin, Harbin, Huludao, Hulunbuir, Heihe, Jiamusi, Jilin, Jinan (bordering Cwa), Linyi (bordering Cwa), Mudanjiang, Qinhuangdao, Qiqihar, Shenyang, Shigatse, Siping, Songyuan, Suihua, Tangshan, Tieling, Ulanhot, Ürümqi (bordering BSk), Xi'an (bordering Cwa), Yanji, Yantai, Yichun
Georgia: Akhaltsikhe
India: Badrinath (bordering Cfb), Dras (Dsb bordering Dfb)
Iran: Abali (Dsb), Arak (Dsa bordering BSk), Ardabil (bordering BSk), Hamadan (Dsa bordering BSk), Saqqez (Dsa), Sardasht (Dsa), Tabriz (Dsa bordering BSk), Urmia (bordering BSk)
Japan: Aomori, Asahikawa, Hakodate, Kushiro, Morioka, Mount Aso, Mutsu, Nagano (bordering Cfa), Nichinan (bordering Cfa), Obihiro, Sapporo, Tendō , Tome (bordering Cfa), Yamagata (bordering Cfa), Yokote
Kazakhstan: Aktobe, Almaty, Arys (Dsa bordering BSk), Astana, Karaganda, Oskemen, Pavlodar, Semey, Shymkent (Dsa), Taraz (Dsa bordering BSk)
Kyrgyzstan: Bishkek (Dsa), Jalal-Abad (Dsa), Karakol, Osh (Dsa)
Mongolia: Baruunturuun (bordering Dwc), Darkhan, Kharkhorin, Sükhbaatar
North Korea: Chongjin, Haeju, Hamhung, Hoeryong, Kaesong, Kimchaek, Nampo, Pyongyang, Rason, Sinuiju, Tanchon, Wonsan
Russia: Abakan (bordering BSk), Barnaul, Birobidzhan, Blagoveshchensk, Chelyabinsk, Chita (bordering Dwc), Gorno-Altaysk, Irkutsk, Kemerovo, Khabarovsk, Krasnoyarsk, Komsomolsk-on-Amur, Kurilsk (bordering Dfc), Kurgan, Nakhodka, Novosibirsk, Omsk, Tomsk, Tyumen, Vladivostok, Yekaterinburg, Yuzhno-Kurilsk, Yuzhno-Sakhalinsk
South Korea: Bucheon, Cheonan, Cheongju, Chuncheon, Chupungnyeong, Daejeon, Goyang, Incheon, Pyeongchang, Seongnam, Seoul, Suwon, Wonju, Yeongcheon (bordering Cfa), Yongin
South Ossetia: Tskhinvali (disputed with Georgia)
Tajikistan: Isfara (Dsa), Konibodom (Dsa), Panjakent (Dsa), Roghun (Dsb)
Turkey: Ağrı (Dsb), Ardahan, Bitlis (Dsa), Çankırı (bordering BSk/Cfa), Çorum (bordering Cfa/Cfb), Erzurum, Hakkâri (Dsa), Kars, Kayseri (Dsa), Muş (Dsa), Sivas (Dsb), Van (Dsa)
Uzbekistan: Chirchiq (Dsa bordering Csa)
Europe
Albania: Pogradec
Andorra: Canillo, El Pas de la Casa (bordering Dfc)
Austria: Baden bei Wien (bordering Cfb) , Innsbruck, Klagenfurt, Klösterle, Neukirchen bei Lambach , Villach, Wiener Neustadt
Belarus: Barysaw, Brest, Gomel, Grodno, Minsk, Pinsk, Vitebsk
Bosnia and Herzegovina: Goražde, Istočno Sarajevo, Livno
Bulgaria: Pazardzhik (bordering Cfa), Pernik, Pleven, Ruse (bordering Cfa), Smolyan (Dsb), Sofia (bordering Cfb), Veliko Tarnovo, Vratsa
Croatia: Gospić
Czech Republic: Brno, České Budějovice, Liberec, Olomouc , Ostrava, Pardubice (bordering Cfb), Plzeň
Estonia: Hiiumaa, Pärnu, Saaremaa, Tallinn, Tartu
Finland: Åland, Hämeenlinna, Helsinki, Kouvola, Kuopio (bordering Dfc), Lahti (bordering Dfc), Lappeenranta, Pori, Tampere (bordering Dfc), Turku
France: Chamonix, Mouthe, Saint-Véran (bordering Dfc)
Germany: Ansbach, Augsburg (bordering Cfb), Bayreuth, Garmisch-Partenkirchen, Görlitz (bordering Cfb), Ingolstadt , Kempten, Regensburg, Sigmaringen, Weiden in der Oberpfalz
Greece: Aetomilitsa, Kato Vermio, Samarina
Hungary: Debrecen, Kecskemét (bordering Cfa/Cfb), Miskolc, Nyíregyháza, Szeged (bordering Cfa/Cfb), Szombathely (bordering Cfb), Szolnok (bordering Cfa)
Italy: Belluno (bordering Cfb), Bruneck, Cortina d'Ampezzo, Rhêmes-Notre-Dame , Rocca di Mezzo (bordering Cfb), Toblach
Kazakhstan: Oral (bordering BSk)
Kosovo: Pristina (disputed with Serbia)
Latvia: Daugavpils, Jelgava, Liepāja, Riga
Liechtenstein: Schaan (bordering Cfb)
Lithuania: Kaunas, Klaipėda, Šiauliai, Vilnius
Moldova: Bălți, Briceni, Chișinău, Comrat
Montenegro: Pljevlja, Žabljak
North Macedonia: Berovo , Bitola
Norway: Bodø (bordering Dfc), Drammen, Hamar, Lillehammer (bordering Dfc), Oslo, Sarpsborg, Steinkjer, Trondheim
Poland: Białystok, Bydgoszcz, Gdańsk, Gorzów Wielkopolski (bordering Cfb), Katowice, Kielce, Kraków, Łódź, Lublin, Olsztyn, Opole, Poznań (bordering Cfb), Rzeszów, Suwałki, Toruń, Warsaw, Zielona Góra (bordering Cfb)
Romania: Brașov, Bucharest, Cluj-Napoca, Craiova, Galați, Iași, Miercurea Ciuc, Oradea, Ploiești, Sibiu
Russia: Belgorod, Bryansk, Cheboksary, Cherkessk, Elista, Grozny, Izhevsk, Kaliningrad, Kazan, Kirov, Kursk, Magas, Moscow, Nalchik, Nizhny Novgorod, Orenburg, Pskov, Perm, Penza, Petrozavodsk (bordering Dfc), Rostov-on-Don, Ryazan, Saint Petersburg, Samara, Saransk, Saratov, Smolensk, Stavropol, Taganrog, Tula, Tver, Ufa, Ulyanovsk, Vladikavkaz, Vladimir, Volgograd (bordering BSk), Voronezh, Yaroslavl, Yeysk, Yoshkar-Ola
Serbia: Nova Varoš, Subotica, Zaječar
Slovakia: Banská Bystrica, Košice, Pezinok , Nitra, Prešov, Trnava, Žilina
Slovenia: Lendava (bordering Cfb), Murska Sobota
Spain: Puerto de Navacerrada (Dsb bordering Csb)
Sweden: Falun, Gävle, Härnösand, Jönköping, Kalmar (bordering Cfb), Karlstad, Linköping, Norrköping, Örebro, Stockholm, Sundsvall, Uppsala, Västerås, Visby (bordering Cfb)
Switzerland: Gstaad, La Brévine (bordering Dfc), La Chaux-de-Fonds, Poschiavo
Transnistria: Tiraspol (disputed with Moldova)
Ukraine: Chernihiv, Dnipro, Donetsk (occupied by Russia), Izmail (bordering Cfa), Kharkiv, Kherson, Khmelnytskyi, Kryvyi Rih, Kyiv, Luhansk (occupied by Russia), Lviv, Mariupol (occupied by Russia), Mykolaiv, Odesa (bordering BSk/Cfa), Poltava, Uzhhorod, Vinnytsia, Zaporizhzhia, Zhytomyr
North America
Canada
Brampton
Brandon, MB
Calgary
Cape Sable Island, NS (bordering Dfc)
Castlegar, BC (Dsb bordering Dfb)
Charlottetown
Edmonton
Fredericton
Gander, NL
Grande Prairie, AB (bordering Dfc)
Goose Bay, NL (bordering Dfc)
Halifax
Hamilton, ON
Kelowna, BC
Kitchener, ON
Laval, QC
London, ON
Longueuil, QC
Marystown, NL
Mississauga
Moncton
Montreal
Moose Factory
Murdochville, QC
Ottawa
Penticton, BC (bordering BSk)
Prince Albert, SK
Prince George, BC
Quebec City
Red Deer, AL
Regina (bordering BSk)
Rimouski, QC
Sable Island (bordering Cfb)
Saint John, NB
Saguenay, QC
Sherbrooke, QC
St. John's, NL
Saskatoon (bordering BSk)
Stewart, BC
Sudbury
Sydney, NS
Thunder Bay
Timmins, ON
Toronto
Trois-Rivières, QC
Windsor
Winnipeg
Yarmouth, NS
United States
Midwest
Aberdeen, SD
Akron, OH
Ann Arbor, MI
Athens, OH
Aurora, IL
Bismarck, ND
Bloomington, IN
Cahokia Heights, IL (bordering Cfa)
Canton, OH
Cedar Rapids, IA
Champaign, IL
Chicago
Chillicothe, OH (bordering Cfa)
Cincinnati (downtown is Cfa)
Cleveland
Columbia, MO
Columbus, OH
Davenport, IA
Dayton
Decatur, IL
Des Moines
Detroit
Duluth, MN
Eau Claire, WI
Fargo, ND
Flint, MI
Fort Wayne, IN
Grand Rapids, MI
Green Bay, WI
Indianapolis
Kalamazoo, MI
Kansas City (bordering Cfa)
Lansing, MI
Lawrence, KS
Lima, OH
Lincoln, NE
Madison, WI
Marietta, OH (bordering Cfa)
Marquette, MI
Milwaukee
Minneapolis
North Platte, NE
Omaha
Oshkosh, WI
Peoria
Rapid City, SD (bordering BSk)
Rochester, MN
Rockford, IL
Saginaw, MI
Sioux Falls, SD
South Bend, IN
Springfield, IL
St. Charles, MO
St. Cloud, MN
St. Joseph, MO
Toledo, OH
Topeka
Youngstown, OH
Northeast
Albany
Allentown, PA
Altoona, PA
Augusta, ME
Bangor, ME
Binghamton, NY
Boston
Bridgeport, CT (bordering Cfa)
Buffalo
Burlington, VT
Cape Cod (Outer Cape and Chatham are Cfb and Woods Hole is Cfa)
Concord, NH
Erie, PA
Harrisburg (downtown and riverfront (including City Island) are Cfa)
Hartford
Ithaca, NY
Johnstown, PA
Lancaster, PA (bordering Cfa)
Lansdale, PA (bordering Cfa)
Lawrence, MA
Lewiston, ME
Lowell, MA
Manchester, NH
Matinicus Isle, ME
Morristown, NJ
Nashua, NH
New Bedford, MA
New Haven
New London, CT (bordering Cfa)
Newport, RI (bordering Cfa/Cfb)
Paterson, NJ (bordering Cfa)
Pittsburgh
Plainfield, NJ (bordering Cfa)
Plymouth, MA (bordering Cfa)
Portland, ME
Portsmouth, NH
Poughkeepsie, NY
Princeton, NJ (bordering Cfa)
Providence, RI
Reading, PA
Riverhead, NY (north side facing Long Island Sound is Cfa)
Rochester
Scranton, PA
Somerville, NJ (bordering Cfa)
Southampton, NY
Springfield, MA
State College, PA
Syracuse
Utica, NY
West Chester, PA (bordering Cfa)
White Plains, NY
Williamsport, PA
Worcester, MA
York, PA (bordering Cfa)
South
Beech Mountain, NC
Boone, NC (bordering Cfb)
Buckhannon, WV (bordering Cfa)
Cincinnati/Northern Kentucky International Airport (bordering Cfa)
Clifton Forge, VA (bordering Cfa)
Cumberland, MD
Ebright Azimuth, DE (bordering Cfa)
Elkins, WV
Emmitsburg, MD (bordering Cfa)
Fairmont, WV (bordering Cfa)
Hancock, MD
Kuwohi, NC/TN
Mount Mitchell, NC
Oakland, MD
Parkton, MD
Romney, WV (bordering Cfa)
Wheeling, WV
White Sulphur Springs, WV
Winchester, VA (bordering Cfa)
West
Alturas, CA (Dsb bordering BSk/Csb)
Anaconda, MT (bordering BSk)
Aspen, CO
Billings, MT (bordering BSk)
Bozeman, MT
Cambridge, ID (Dsa)
Cheyenne, WY (bordering BSk)
Coeur D'Alene, ID (Dsb bordering Csb)
Fairbanks (bordering Dfc)
Flagstaff, AZ (bordering BSk)
Government Camp, OR (Dsb bordering Csb/Csc/Dsc)
Idaho Falls (bordering BSk)
Jackson, WY (bordering Dfc)
Joseph, OR
Juneau (bordering Dfc)
Kalispell, MT
Klamath Falls, OR (Dsb)
Kodiak, AK (bordering Cfb/Cfc/Dfc)
La Grande, OR (Dsb bordering Csb)
Logan, UT (Dsa)
Los Alamos, NM (bordering BSk)
Loveland, CO (bordering BSk)
Mammoth Lakes, CA (Dsb)
Missoula, MT
Moscow, ID (Dsb)
Mountain City, NV (bordering BSk)
Nenana, AK (bordering Dfc)
Ogden, UT (Dsa)
Orofino, ID (Dsb bordering Dsa)
Park City, UT
Petersburg, AK
Pocatello, ID
Rexburg, ID
Salt Lake City (Dsa bordering BSk/Csa)
Santa Fe, NM (bordering BSk)
Skagway, AK (Dsb bordering Dsc)
South Lake Tahoe, CA (Dsb bordering Csb)
Spokane, WA (Dsb bordering Csa/Csb/Dsa)
Steamboat Springs, CO
Tahoe City, CA (Dsb)
Taos, NM (bordering BSk)
Telluride, CO
Tusayan, AZ (bordering BSk)
Winthrop, WA (Dsb)
Oceania
Australia
Falls Creek, Victoria (bordering Dfc)
Kiandra, New South Wales (bordering Cfb)
Mount Buller, Victoria (bordering Cfb/Cfc/Dfc)
Perisher Valley, New South Wales (bordering Cfb/Cfc/Dfc)
South America
Argentina
Alto Río Senguer, Chubut Province (Dsb bordering BSk/Csb/Csc/Dsc)
Las Leñas, Mendoza Province (Dsb)
Puente del Inca, Mendoza Province (Dsb bordering Csb)
| Physical sciences | Climates | Earth science |
249773 | https://en.wikipedia.org/wiki/Killifish | Killifish | A killifish is any of various oviparous (egg-laying) cyprinodontiform fish, including families Aplocheilidae, Cyprinodontidae, Fundulidae, Nothobranchiidae, Profundulidae, Aphaniidae and Valenciidae. All together, there are 1,270 species of killifish, the biggest family being Rivulidae, containing more than 320 species. As an adaptation to living in ephemeral waters, the eggs of most killifish can survive periods of partial dehydration. Many of the species rely on such a diapause, since the eggs would not survive more than a few weeks if entirely submerged in water. The adults of some species, such as Kryptolebias marmoratus, can additionally survive out of the water for several weeks. Most killifish are small, measuring from , with the largest species growing to just under .
The word killifish is of uncertain origin, but is likely to have come from the Dutch kil for a kill (small stream). Although killifish is sometimes used as an English equivalent to the taxonomical term Cyprinodontidae, some species belonging to that family have their own common names, such as the pupfish and the mummichog.
Range and habitat
Killifish are found mainly in fresh or brackish waters in the Americas, as far south as Argentina and as far north as southern Ontario and even Newfoundland and Labrador. There are also species in southern Europe, in much of Africa as far south as KwaZulu-Natal, South Africa, in the Middle East and Asia (as far east as Vietnam), and on several Indian Ocean islands.
The majority of killifish are found in permanent streams, rivers, and lakes, and live between two and three years. Such killifish are common in the Americas (Cyprinodon, Fundulus and Rivulus) as well as in Africa and Asia (including Aphyosemion, Aplocheilus, Epiplatys, Fundulopanchax and Lacustricola) and southern Europe (Aphanius). Some of these habitats can be rather extreme; the only natural habitat of the Devils Hole pupfish (Cyprinodon diabolis) is Devils Hole: a cavern at least deep, branching out from a small opening at the surface, approximately by wide.
Some specialized forms live in temporary ponds and flood plains, and typically have a much shorter lifespan. Such species, known as "annuals", live no longer than nine months, and are used as models for studies on aging. Examples include the African genus Nothobranchius and South American genera ranging from the cold water Austrolebias of Argentina and Uruguay to the more tropical Gnatholebias, Pterolebias, Simpsonichthys and Terranatos.
Territorial behaviour
A small number of species will shoal while most are territorial to varying degrees. Populations can be dense and territories can shift quickly, especially for species of the extreme shallows (a few centimetres of water). Many species exist as passive tribes in small streams where dominant males will defend a territory while allowing females and immature males to pass through the area. In the aquarium, territorial behaviour is different for every grouping, and will even vary by individuals. In a large enough aquarium, most species can live in groups as long as there are more than three males.
Diet
Killifish feed primarily on aquatic arthropods such as insect (mosquito) larvae, aquatic crustaceans and worms. Some species of Orestias from Lake Titicaca are planktonic filter feeders. Others, such as Cynolebias and Megalebias species and Nothobranchius ocellatus are predatory and feed mainly on other fish. The American Flagfish (Jordanella floridae) feeds heavily on algae and other plant matter as well as aquatic invertebrates. Nothobranchius furzeri needs much food because it grows quickly, so when food supplied is inadequate, bigger fish will eat the smaller fish.
In lifespan research
Some strains have a lifespan as short as several months and can thus serve as a model for biogerontological studies. The African turquoise killifish (Nothobranchius furzeri) is the shortest-living vertebrate that can be bred in captivity, having a lifespan of between three and nine months. Sexual maturation occurs within 3–4 weeks, with fecundity peaking in 8–10 weeks.
Nothobranchius furzeri shows no signs of telomere shortening, reduced telomerase activity, or replicative senescence with age, despite its short lifespan. Nonetheless, lipofuscin accumulates in the brain and liver (associated with age-related neurodegeneration), and there is an increased risk of cancer with age. Calorie restriction reduces these age-related disease conditions. Resveratrol has been shown to increase the mean (56%) and maximum life span (59%) of Nothobranchius furzeri, but resveratrol has not been shown to have this effect in mammals.
Transferring the gut microbiota from young killifish into middle-aged killifish significantly extends the lifespans of the middle-aged killifish.
Transgenic strains have been made, and precise genome editing was achieved in Nothobranchius furzeri using a draft genome and the CRISPR/Cas9 system. By targeting multiple genes, including telomerase, the killifish can now be used as an attractive vertebrate model organism for aging and diseases (such as Dyskeratosis congenita). Sequencing the whole killifish genome indicated modification to the IGF-1 receptor gene.
As pets
Many killifish are lavishly coloured and most species are easy to keep and breed in an aquarium. Specimens can be obtained from specialist societies and associations. Striped panchax (also known as the Golden Wonder killifish) are commonly found in pet shops, but caution must be exercised when considering tank mates, since the mouth of the Striped panchax is as wide as the head, and much smaller fish will be eaten. Flagfish, native to south Florida, is another species of killifish commonly found in pet stores. They are useful in aquariums for algae control. The golden topminnow (Fundulus chrysotus) is also native to the United States and often available at auction, but care must be taken with these fish to stop them from jumping out of the tank. A firm cover and a layer of floating plants is best when keeping these fish.
Behaviour-altering infection
Normally, killifish avoid near-surface water to reduce the danger of being eaten by predators. However, when infected with a type of fluke the fish swim near the surface, and sometimes even swim upside down, exposing their camouflaged bellies. The fluke completes its lifecycle in the digestive tract of birds.
Evolved resistance to extreme levels of toxicity
The large populations of killifish and the genetic diversity of the species have enabled it to evolve and survive in areas where other species have died out, including Superfund sites. Over a few dozen generations of killifish in a relatively short period of time (50–60 years), killifish have evolved resistance against levels of dioxins, PCBs, mercury, and other industrial chemicals up to 8,000 times higher than the previously estimated lethal dose. Sequencing the genomes of the adapted individuals showed a common set of mutations among the pollution-tolerant fish, many of which help to deactivate or turn off a molecular pathway responsible for a large part of the cellular damage caused by the chemicals.
Killifish were found to fare relatively well in the wake of the Deepwater Horizon oil spill.
| Biology and health sciences | Acanthomorpha | Animals |
3484275 | https://en.wikipedia.org/wiki/Dean%E2%80%93Stark%20apparatus | Dean–Stark apparatus | The Marcusson apparatus, Dean-Stark apparatus, Dean–Stark receiver, distilling trap, or Dean–Stark Head is a piece of laboratory glassware used in synthetic chemistry to collect water (or occasionally other liquid) from a reactor. It is used in combination with a reflux condenser and a distillation flask for the separation of water from liquids. This may be a continuous removal of the water that is produced during a chemical reaction performed at reflux temperature, such as in esterification reactions. The original setup by Julius Marcusson (invented in 1905) was refined by the American chemists Ernest Woodward Dean (1888–1959) and David Dewey Stark (1893–1979) in 1920 for determination of the water content in petroleum.
Function
Two types of Dean–Stark traps exist – one for use with solvents with a density less than that of water and another for use with solvents with a density greater than that of water.
The Dean–Stark apparatus typically consists of a vertical cylindrical glass tube, often with a volumetric graduation along its full length and a precision stopcock at its lower end, very much like a burette. Traps designed to remove or measure very small amounts of water may be closed, with no tap. The lower end of a reflux condenser fits into the top of the cylinder. Immediately below the joint between the condenser and the cylinder is a sloping side-arm that joins the cylinder to a reaction flask. The lower end of the side-arm turns sharply downward, so that the side-arm is connected to the reaction flask by a vertical tube.
The reaction flask is heated. Boiling chips within it assist with the calm formation of bubbles of vapor containing the reaction solvent and the component to be removed. This vapor travels out of reaction flask up into the condenser where water being circulated around it causes it to cool and drip into the distilling trap. Here, the immiscible liquids separate into layers (water below and solvent above it). When their combined volume reaches the level of the side-arm, the upper, less-dense layer will begin to flow back to the reactor while the water layer will remain in the trap. The trap will eventually reach capacity when the level of the water in it reaches the level of the side-arm. At this point, the trap must be drained into the receiving flask. The process of evaporation, condensation and collection may be continued until it ceases to produce additional amounts of water.
More rarely encountered is the model for solvents with a density greater than water. This type has a tube at the bottom of the collection vessel to allow the organic solvent at the bottom to flow back into the reaction vessel. The water generated during the reaction floats on top of the organic phase. The entire volume must be drained to remove the water and the organic layer can be separated to be returned to the system.
This piece of equipment is usually used in azeotropic distillations. A common example is the removal of water generated during a reaction in boiling toluene. An azeotropic mixture of toluene and water distills out of the reaction, but only the toluene (density 0.865g/ml) returns, since it floats on top of the water (density 0.998g/ml), which collects in the trap. Some high-boiling liquids that have an azeotrope with water can be dried by adding toluene or another azeotrope-breaking solvent to allow the extraction of water. The Dean–Stark method is commonly used to measure moisture content of items such as bread in the food industry.
This equipment can be used in cases other than simple removal of water. One example is the esterification of butanol with acetic acid catalyzed by sulfuric acid. The vapor contains 63% ester, 29% water and 8% alcohol at reflux temperature and the organic layer in the trap contains 86% ester, 11% alcohol and 3% water which is reintroduced. The water layer is 97% pure.
Another example is the esterification of benzoic acid and n-butanol where the ester product is trapped and the butanol, immiscible with the water, flows back into the reactor. Removing water in the course of these esterifications shifts the chemical equilibrium in favor of ester formation, in accordance with Le Chatelier's principle.
| Physical sciences | Phase separations | Chemistry |
3486181 | https://en.wikipedia.org/wiki/Rigid%20airship | Rigid airship | A rigid airship is a type of airship (or dirigible) in which the envelope is supported by an internal framework rather than by being kept in shape by the pressure of the lifting gas within the envelope, as in blimps (also called pressure airships) and semi-rigid airships. Rigid airships are often commonly called Zeppelins, though this technically refers only to airships built by the Luftschiffbau Zeppelin company.
In 1900, Count Ferdinand von Zeppelin successfully performed the maiden flight of his first airship; further models quickly followed. Prior to the First World War, Germany was a world leader in the field, largely attributable to the work of von Zeppelin and his Luftschiffbau Zeppelin company. During the conflict, rigid airships were tasked with various military duties, which included their participation in Germany's strategic bombing campaign. Numerous rigid airships were produced and employed with relative commercial success between the 1900s and the late 1930s. The heyday of the rigid airship was abruptly ended by the destruction of the Hindenburg by fire on 6 May 1937. The disaster not only destroyed the biggest zeppelin in the world, but the film caused considerable reputation damage to rigid airships in general. Several nations had ended military rigid airship programs after serious accidents earlier in the decade, but widespread public safety concerns in the wake of the Hindenburg disaster led several nations to permanently ground their existing rigid airships and scrap them in subsequent years.
Construction and operation
Rigid airships consist of a structural framework usually covered in doped fabric containing a number of gasbags or cells containing a lifting gas. In the majority of airships constructed before the Second World War, highly flammable hydrogen was used for this purpose, resulting in many airships such as the British R101 and the German Hindenburg being lost in catastrophic fires. The inert gas helium was used by American airships in the 1920s and 1930s; it is also used in all modern airships.
Airships rely on the difference in density between the lifting gas and the surrounding air to stay aloft. Typically airships start a flight with their gasbags inflated to about 95% capacity: as the airship gains height the lifting gas expands as the surrounding atmospheric pressure reduces. As the surrounding atmospheric pressure decreases, the lifting gas expands, displacing ambient air. When the entire envelope is filled with expanded lifting gas, the aircraft is at its pressure height, which is generally the maximum operational ceiling. At this point, excess expanding gas must either be vented or the airship must descend so that the lifting gas can contract and ambient air brought back into the hull.
Airships can also generate a certain amount of aerodynamic lift by using their elevators to fly in a nose-up attitude. Similarly, by flying nose-down, down-force can be generated: this may be done to prevent the airship rising above its pressure height.
History
Early history
By 1874, several people had conceived of a rigid dirigible (in contrast to non-rigid powered airships which had been flying since 1852). The Frenchman Joseph Spiess had patented a rigid airship design in 1873 but failed to get funding. Another such individual was the German Count Ferdinand von Zeppelin, who had outlined his thoughts of a rigid airship in diary entries from 25 March 1874 through to 1890 when he resigned from the military. David Schwarz had thought about building an airship in the 1880s and had probably started design work in 1891: by 1892, he had started construction.
However, Schwarz's all-aluminium airship would not perform any test flights until after his death in 1897. Schwarz had secured help in its construction from the industrialist Carl Berg and the Prussian Airship Battalion; there was an exclusive contract in place between Schwarz and Berg, thus Count Zeppelin was obliged to reach a legal agreement with Schwarz's heirs to obtain aluminium from Carl Berg, although the two men's designs were different and independent from each other: the Schwarz design lacked the separate internal gasbags that characterise rigid airships. Using Berg's aluminium, von Zeppelin was able to start building his first airship, the LZ 1, in 1899.
First practical rigid airships
During July 1900, Ferdinand von Zeppelin completed LZ 1. Constructed in a floating shed on Lake Constance, it was 128.02 m (420 ft) long, 11.73 m (38 ft 6 in) in diameter with a volume of 11,298 m3 (399,000 ft 3) and was powered by a pair of 11 kW (14 hp) Daimler engines. The first flight, lasting 20 minutes, was made on 2 July, but ended with the airship being damaged. After repairs and modifications, two further flights were conducted in October 1900. However, these initial experiments failed to attract any investors, and Count Zeppelin did not complete his next design, LZ 2, until 1906. This performed only a single flight on 17 January 1906, during which both engines failed and the zeppelin was compelled to conduct a forced landing in the Allgäu mountains; it was subsequently damaged beyond repair by a storm. Undeterred, another zeppelin with a largely similar design, the LZ 3, was quickly completed and put into flight.
LZ 3 proved to have performed sufficiently to interest the German Army, who opted to purchase and operate it as the Z I until 1913. Even so, the German Army observed that they required an airship that would be capable of flying for 24 hours. As this was beyond the capability of LZ 3, it was decided to design and construct a larger craft, LZ 4. This was 136 m (446 ft) long, 12.95 m (42 ft 6 in) in diameter and powered by two Daimler engines delivering a total of 156 kW (210 hp). LZ 4 first flew on 20 June 1908, and on 1 July made a spectacular 12 hour cross-country flight during which it was flown over Switzerland to Zürich and then back to Lake Constance. The 24-hour trial was started on 4 August, but was interrupted by the failure of one of the engines. It was moored near Echterdingen in order to make repairs but a storm arose, causing it to break away from its moorings, after which it was blown into some trees and caught fire. The disaster took place in front of an estimated 40 to 50 thousand spectators, and produced an extraordinary wave of nationalistic support for von Zeppelin's work. Unsolicited donations from the public poured in: enough had been received within 24 hours to rebuild the airship, and the eventual total was over 6 million marks were donated, finally giving Count Zeppelin a sound financial base for his experiments.
Seven zeppelins were operated by DELAG, the first airline in the world. DELAG was founded at the suggestion of Alfred Colsman, the business manager of Zeppelin Luftschiffbau, seeking to capitalise on the German public's enthusiastic interest in the zeppelin by permitting them onboard passenger-carrying airships as a commercial venture; von Zeppelin distanced himself from this commercialisation, reportedly regarding such efforts to have been a vulgar tradesman's enterprise. Commencing such flights in 1910, DELAG was initially limited to offering pleasure cruises in the vicinity of the existing zeppelin bases.
DELAG soon received more capable zeppelins, such as the LZ 10 Schwaben, which would carry a total of 1,553 paying passengers during its career, which involved not only pleasure flights but a number of long-distance flights to destinations such as Frankfurt, Düsseldorf, and Berlin. The company's airships were also used by the Imperial German Navy for crew training, with the Navy crews operating passenger flights. By July 1914, one month prior to the start of the First World War, DELAG's Zeppelins had transported a total of 34,028 passengers on 1,588 commercial flights; over these trips, the fleet had accumulated 172,535 kilometres across 3,176 hours of flight. Commercial operations came to an abrupt end in Germany due to the outbreak of the First World War, after which DELAG's airships were taken over by the German Army for wartime service.
During 1911, the first rigid airship produced by the German Schütte-Lanz company was flown. Designed by the naval architect Johann Schütte, the Schütte-Lanz introduced a number of technical innovations. The shape of the hull was more streamlined than the early Zeppelin craft, the hulls of which were cylindrical for most of their length, simplifying construction at the expense of aerodynamic efficiency. Other Schütte-Lanz innovations included the use of an axial cable running the length of the airship to reduce additional stressing caused by the partial deflation of a single gasbag, the introduction of venting tubes to carry any hydrogen vented to the top of the ship and simplified cruciform tail surfaces.
The British Royal Navy took an early interest in rigid airships and ordered His Majesty's Airship No. 1 in 1909 from Vickers Limited at Barrow-in-Furness. It was 512 ft (156.06 m) long with two Wolseley engines. It was completed in 1911 but broke in two before its first flight and was scrapped. This caused a temporary halt to British airship development, but in 1913 an order was placed for HMA No. 9r. Due to various factors, including difficulties in acquiring the necessary materials, it was not completed until April 1917.
France's only rigid airship was designed by Alsatian engineer Joseph Spiess and constructed by Société Zodiac at the Aérodrome de Saint-Cyr-l'École. It had a framework of hollow wooden spars braced with wire, and was given the name Zodiac XII but had the name SPIESS painted along the side of the envelope. It was 113 m (370 ft 9 in) long, with a diameter of 13.5 m (44 ft 3 in) and was powered by a single Chenu 200 hp engine that drove two propellers. It first flew on 13 April 1913, but it became clear that it was underpowered and required more lift, so it was lengthened to 140 m (459 ft 4 in) to accommodate three more gas cells and a second engine was added. Spiess then presented the airship to the French government as a gift. After further trials it was not accepted by the French military, because their view was that smaller non-rigid types would be more effective. The Spiess airship seems to have been broken-up in 1914.
First World War
During the First World War, the Zeppelin company constructed a total of 95 military airships. These were operated by both the German Navy and the Army. German military airship stations had been established before the conflict and on September 2–3, 1914, the Zeppelin LZ 17 dropped three 200 lb bombs on Antwerp in Belgium. In 1915, a bombing campaign against England using airships was initiated, the first raid taking place on 19 January 1915 when two airships dropped bombs on Norfolk. On 31 May 1915 the first bombs fell on London. Raids continued throughout 1915 and continued into 1916. On the night of September 2–3, 1916 the first German airship was shot down over English soil by Lt. Leefe Robinson flying a BE 2c. This and subsequent successes by Britain’s defences led to the development of new Zeppelin designs capable of operating at greater altitudes, but even when these came into service the Germans only carried out a small number of airship raids on Britain during the rest of the war, carrying on the campaign using aeroplanes and reserving their airships for their primary duty of naval patrols over the North Sea and the Baltic. The last casualties occurred on 12 April 1918.
The first British airship to be completed during the war was No. 9r, which was first flown at the end of 1916 and was used for experimental and training purposes. By then, the war against U-boats was at its height and 9r was quickly followed by four airships of the 23 Class, two R23X Class and two R31 Class, the last being based on the Schütte-Lanz principle of wooden construction, and remain the largest mobile wooden structures ever built. The only significant combat success of these airships, aside from their deterrent effect, was assistance in the destruction of SM UB-115 by R29 in September 1918.
1919–1939
By the end of the conflict, two British airships of the R33 Class were nearing completion. R33 became a civilian airship, finishing her career doing experimental work. The R34 became the first aircraft to complete a return Atlantic crossing in July 1919 but was severely damaged in January 1921 and was subsequently scrapped. R.35, a unique admiralty design, was almost finished when work was stopped in early 1919. R36 and R.37 were stretched R.35s. R.36 was completed after the war as a civilian airship registered as G-FAAF. R.36 had two engines from the German L71. Modifications for passenger service involved installing a 131 foot long combined control and passenger gondola to accommodate 50 passengers. R.36 suffered a structural failure of one horizontal and one vertical fin. It was repaired and served to aid the police in traffic control for the Ascot race in 1921. R.36 was damaged in a mooring accident in 1921, and while repaired R.36 never flew again. Retained for possible use as a commercial airship R.36 was broken up in 1926.
Four airships of the R38 Class were started but only one completed: it was sold to the US Navy and renamed ZR-2. In June 1921 it broke up in the air over Kingston-upon-Hull before it could be delivered, killing 44 of its Anglo-American crew. The last airship that had been ordered amid the First World War was the R80; it was completed in 1920 but was tested to destruction in the following year after it was found to have no commercial use.
After the end of World War I, Luftschiffbau Zeppelin resumed building and operating civilian airships. Under the terms of the Treaty of Versailles, Germany was prohibited from building airships with a capacity in excess of , greatly limiting the company's scope. However, a pair of small passenger airships, LZ 120 Bodensee and a sister ship LZ 121 Nordstern were built, intended for use between Berlin and Friedrichshafen. They were subsequently confiscated and handed over to Italy and France as war reparations in place of wartime zeppelins which had been sabotaged by their crews in 1919. The Zeppelin company was saved from extinction by an order for an airship, the USS Los Angeles, being placed by the US Navy; this airship conducted its first flight on 27 August 1924. The Goodyear-Zeppelin partnership would continue up until the outbreak of the Second World War.
In 1924, the British Government initiated the Imperial Airship Scheme, a plan to launch airship routes throughout the British Empire. This involved the construction of two large airships, the R100 and R101, paid for by the government. The R100 was privately built by Vickers-Armstrongs using existing commercial practices, with a design team led by Barnes Wallis, who had previously co-designed the R80. After her first flight in December 1929, R100 made a successful round trip to Quebec in Canada in July and August the following year. The competing R101 was designed and built by the Air Ministry and was supposed to encourage new approaches. R101 was severely overweight, largely due to the decision to use diesel engines to reduce fire risk, and it was decided to lengthen the airship's hull to increase lift. In October 1930, R101 set off to Karachi on its first overseas flight but crashed in northern France in bad weather killing 48 of the 54 people on board, including the Secretary of State for Air and most of the design team. Following this disaster, the R100 was grounded and was finally scrapped in November 1931, marking the end of British interest in rigid airships.
During 1925, the Versailles restrictions were relaxed by the Allies, enabling Dr Hugo Eckener, the chairman of Zeppelin Luftschiffbau, to pursue his vision of developing a zeppelin suitable for launching an intercontinental air passenger service. The sum of 2.5 million Reichsmarks (ℛℳ, the equivalent of US$600,000 at the time, or $ million in 2018 dollars), was raised via public subscription, while the German government also granted over ℛℳ 1 million ($ million) for the project. Accordingly, Zeppelin Lufftschiffbau began construction of the first of a new generation of airships, the LZ 127 Graf Zeppelin. On 18 September 1928, the completed airship flew for the first time. Shortly thereafter, DELAG commenced operations with the Graf Zeppelin, being enabled to launch regular, nonstop, transatlantic flights several years before airplanes would be capable of sufficient range to cross the ocean in either direction without stopping. During 1931, the Graf Zeppelin began offering regular scheduled passenger service between Germany and South America, a route which was continued up until 1937. During its career, Graf Zeppelin crossed the South Atlantic a total of 136 times. The airship also performed numerous record-breaking flights, including a successful circumnavigation of the globe.
The United States rigid airship program was based at Lakehurst Naval Air station, New Jersey. was the first rigid airship constructed in America, and served from 1923 to 1925, when it broke up in mid-air in severe weather, killing 14 members of its crew. was a German airship built for the United States in 1924. The ship was grounded in 1931, due to the Depression, but was not dismantled for over 5 years. A pair of large airships, the Akron and Macon, that both functioned as flying aircraft carriers were procured by the US Navy. However, they were both destroyed in separate accidents. the Akron was flown into the sea in bad weather and broke up, resulting in the deaths of over seventy people, including one of the US Navy's proponents of airships, Rear Admiral William A. Moffett. Macon also ended up in the sea when it flew into heavy weather with unrepaired damage from an earlier incident, but the introduction of life-jackets following the loss of the Akron meant only two people died.
LZ 129 Hindenburg carried passengers, mail and freight on regularly scheduled commercial services from Germany to North and South America. However, such services were brought to an abrupt end by the Hindenburg disaster of 1937. While the Hindenburg's sister ship, the LZ 130 Graf Zeppelin II, was completed, it would only perform thirty European test and government-sponsored flights before being grounded permanently. During 1938, Luftschiffbau Zeppelin was compelled to terminate Zeppelin manufacturing, while all operations of existing airships was ceased within two years. The frames of Graf Zeppelin and Graf Zeppelin II, along with scrap material from the Hindenburg, were subsequently scrapped that same year for their materials, which were used to fulfil wartime demands for fixed-wing military aircraft for the Luftwaffe.
Demise
Following the Hindenburg disaster, the Zeppelin company resolved to use helium in their future passenger airships. However, by this time, Europe was well on the path to the Second World War, and the United States, the only country with substantial helium reserves, refused to sell the necessary gas. Commercial international aviation was limited during the war, so development of new airships was halted. Although several companies, including Goodyear, proposed post-war commercial designs, these were largely to no avail. At an Air Ministry post-war planning session in 1943, a R.104 was proposed to fulfill the Air Ministry Specification C.18/43. Despite the presence of two airship stalwarts, Nevil Shute and Wing Commander T.R. Cave-Browne-Cave the airship was not adopted. The proposed R.104 was described by Lord Beaverbrook as "A pretty face, but no good in the kitchen." The decision was to develop the Bristol Brabazon to meet C.18/43. The Brabazon was a much ballyhooed failure of the post war period. Following the rapid advances in aviation during and after World War II, fixed-wing heavier-than-air aircraft, able to fly much faster than rigid airships, became the favoured method of international air travel.
Modern rigids
The last rigid airships designed and built were built in the 1960s. The AEREON III was constructed in Mercer County, New Jersey in the mid-1960s. It was to utilize the method of "propulsion" developed and demonstrated by Doctor Solomon Andrews in the 1860s as well as a stern-mounted engine. The AEREON III, which had three side-by-side hulls, flipped over during taxi tests and was never repaired. A replacement, the AEREON 26, with a delta configuration, was constructed and flight-tested in the early 1970s. The test program ended due to the expiration of the life time of the drone engine. It was last reported hangared at the Trenton-Robbinsvile Airport in New Jersey. It is not known whether it still exists after almost 50 years.
The Zeppelin company refers to their NT ship as a rigid, but the envelope shape is retained in part by super-pressure of the lifting gas, and so the NT is more correctly classified as semi-rigid.
Aeroscraft was certified airworthy by the FAA in September 2013 and has begun flight testing.
In 2023, the Pathfinder 1, a prototype electric airship by LTA Research, was unveiled. It is the largest modern airship at 124.5 metres long.
| Technology | Types of aircraft | null |
3488015 | https://en.wikipedia.org/wiki/Stunted%20growth | Stunted growth | Stunted growth, also known as stunting or linear growth failure, is defined as impaired growth and development manifested by low height-for-age. It is a manifestation of malnutrition (undernutrition) and can be caused by endogenous factors (such as chronic food insecurity) or exogenous factors (such as parasitic infection). Stunting is largely irreversible if occurring in the first 1000 days from conception to two years of age. The international definition of childhood stunting is a child whose height-for-age value is below -2 standard deviations from the median of the World Health Organization's (WHO) Child Growth Standards. Stunted growth is associated with poverty, maternal under-nutrition and poor health, frequent illness, and/or inappropriate feeding practices and care during the early years of life.
, an estimated 149 million children under 5 years of age are stunted worldwide. More than 85% of the world's stunted children live in Africa and Asia. Once established, stunting and its effects typically become permanent. Stunted children may never regain the height lost as a result of stunting, and most children will never gain the corresponding body weight. One notable contribution to stunted growth is a lack of sanitation, such as public defecation in countries like India.
Health effects
Stunted growth in children has the following public health impacts (apart from the obvious impact of shorter stature of the person affected):
Greater risk for illness and premature death
Delayed cognitive development, poor school performance
Reduced intelligence quotient
Future risk of obesity
Women of shorter stature have a greater risk for complications during child birth due to their smaller pelvis, and are at risk of delivering a baby with low birth weight
Stunted growth can be passed to the next generation, known as the "inter-generational cycle of malnutrition"
The impact of stunting on child development has been established in multiple studies. If a child is stunted at age 2, they tend to have a higher risk of poor cognitive and educational achievement in life, with subsequent socioeconomic and inter-generational consequences. Multi-country studies have also suggested that stunting is associated with reductions in schooling, decreased economic productivity and poverty. Stunted children also display higher risk of developing chronic non-communicable conditions such as diabetes and obesity as adults. If a stunted child undergoes substantial weight gain after age 2, this can lead to obesity. This is believed to be caused by metabolic changes produced by chronic malnutrition that can produce metabolic imbalances if the individual is exposed to excessive or poor quality diets as an adult. This can lead to a higher risk of developing other related non-communicable diseases such as hypertension, coronary heart disease, metabolic syndrome and stroke.
At societal level, stunted individuals do not fulfill their physical and cognitive developmental potential and will not be able to contribute maximally to society. Stunting can therefore limit economic development and productivity, and it has been estimated that it can affect a country's GDP by up to 3%.
Stunting is prevalent in the Global South and has severe consequences including increased risk of infections and death. The global percentage of stunted growth decreased from 33% to 22.3% between 2000 and 2022. The largest drop took place in Asia, from 37.1% in 2000 to 28.2% in 2012 and 22.3% in 2022. Despite global progress, the prevalence of child stunting was greater than 30 percent in 28 countries in 2022 (most of which in sub-Saharan Africa).
Causes
In many publications, the causes for stunting are considered very similar if not the same as the causes for malnutrition in children. However, there are some that contradict this notion. Recent evidence suggests that stunting may not be taken as a synonym of malnutrition, but as the natural condition of human height in non-Westernized societies.
Almost all stunting occurs within the 1,000-day period that spans from conception to a child's second birthday, which constitutes a window of opportunity for growth promotion. The recognition of prenatal factors underlines the inter-generational aspects of growth, and the need for early interventions. The three main causes of stunting in South Asia, and probably in most developing countries, are poor feeding practices, poor maternal nutrition, and poor sanitation. A recent risk assessment analysis for 137 developing countries found that the leading risk factors for stunting were fetal growth restriction (birth weight <10th percentile) followed by unimproved sanitation and diarrhea. It was estimated that 22% of stunting cases were attributable to environmental factors while 14% were attributable to child nutrition. In addition, looking at trends from 1970 to 2012 for 116 countries, women’s education, gender equality and finally quantity and quality of foods available at the country level have been instrumental in reducing stunting rates, while income growth and governance have played facilitating roles.
Feeding practices
Inadequate complementary child feeding and a general lack of vital nutrients beside pure caloric intake is one cause for stunted growth. Children need to be fed diets which meet the minimum requirements in terms of frequency and diversity in order to prevent under-nutrition. Exclusive breastfeeding is recommended for the first six months of life and complementary feeding of nutritious food alongside breastfeeding for children aged six months to 2-years-old. Prolonged exclusive breastfeeding is associated with under-nutrition because breast milk alone is nutritionally insufficient for children over six months old. Breastfeeding for a long time with inadequate complementary feeding leads to growth failure due to insufficient nutrients which are essential for childhood development. The relationship between under-nutrition and prolonged duration of breastfeeding is mostly observed among children from poor households with uneducated parents, as they are more likely to continue breast-feeding without meeting minimum dietary diversity requirement.
Maternal nutrition
Poor maternal nutrition during pregnancy and breastfeeding can lead to stunted growth of children. Proper nutrition for mothers during the prenatal and postnatal period is important for ensuring healthy birth weight and for healthy childhood growth. Prenatal causes of child stunting are associated with maternal undernutrition. Low maternal BMI predisposes the fetus to poor growth leading to intrauterine growth retardation, which is strongly associated with low birth weight and size. Women who are underweight or anemic during pregnancy, are more likely to have stunted children which perpetuates the inter-generational transmission of stunting. Children born with low birth-weight are more at risk of stunting. However, the effect of prenatal under-nutrition can be addressed during the postnatal period through proper child feeding practices.
Maternal under-nutrition increases the risk of stunting at 2 years of age. Based on data from 19 birth cohorts from LMICs, 20% of stunting is attributed to being born small-for-gestational-age (SGA). Further, estimated stunting at 2 years attributed to fetal growth restriction and preterm birth in 2011 was 33% in all developing countries and 41% in South Asia. Restricted pre- and postnatal growth are in turn important determinants of short adult height, increasing the likelihood of the next generation experiencing stunted growth.
Sanitation
There may be a link between children's linear growth and household sanitation practices. The ingestion of high quantities of fecal bacteria by young children through putting soiled fingers or household items in the mouth leads to intestinal infections. This affects children's nutritional status by diminishing appetite, reducing nutrient absorption, and increasing nutrient losses. The diseases recurrent diarrhea and intestinal worm infections (helminthiasis) which are both linked to poor sanitation have been shown to contribute to child stunting.
Research on a global level has found that the proportion of stunting that could be attributed to five or more episodes of diarrhea before two years of age was 25%. Since diarrhea is closely linked with water, sanitation and hygiene (WASH), this is a good indicator for the connection between WASH and stunted growth. To what extent improvements in drinking water safety, toilet use and good handwashing practices contribute to reduce stunting depends on the how bad these practices were prior to interventions.
Environmental enteropathy
The condition termed environmental enteropathy is proposed as an immediate causal factor of childhood stunting. This is an asymptomatic small intestinal disorder characterized by chronic gut inflammation, reduced absorptive surface area, and disruption of intestinal barrier function. This small bowel disorder can be attributed to sustained exposure to intestinal pathogens caused by fecal contamination of food and water. Recent evidence confirmed a causal relationship between stunted growth and environmental enteropathy in children. Several studies are also underway to examine the link between this condition and stunted growth. The exact parthenogenesis of environmental enteropathy causing linear growth failure is unclear, but it is hypothesized that chronic inflammatory state and impaired absorption associated with this condition may inhibit bone growth and affect the linear growth during early years of life.
Diagnosis
Growth stunting is identified by comparing measurements of children's heights to the World Health Organization 2006 growth reference population: children who fall below the fifth percentile of the reference population in height for age are defined as stunted, regardless of the reason. The lower than fifth percentile corresponds to less than two standard deviations of the WHO Child Growth Standards median.
As an indicator of nutritional status, comparisons of children's measurements with growth reference curves may be used differently for populations of children than for individual children. The fact that an individual child falls below the fifth percentile for height for age on a growth reference curve may reflect normal variation in growth within a population: the individual child may be short simply because both parents carried genes for shortness and not because of inadequate nutrition. However, if substantially more than 5% of an identified child population have height for age that is less than the fifth percentile on the reference curve, then the population is said to have a higher-than-expected prevalence of stunting, and malnutrition is generally the first cause considered.
Prevention
Three main things are needed to reduce stunting:
an environment where political commitment can thrive (also called an "enabling environment")
applying several nutritional modifications or changes in a population on a large scale which have a high benefit and a low cost
a strong foundation that can drive change (food security, and a supportive health environment through increasing access to safe water and sanitation).
To prevent stunting, it is not just a matter of providing better nutrition but also access to clean water, improved sanitation (hygienic toilets) and hand washing at critical times (summarized as "WASH"). Without provision of toilets, prevention of tropical intestinal diseases, which may affect almost all children in the developing world and lead to stunting will not be possible.
Studies have looked at ranking the underlying determinants in terms of their potency in reducing child stunting and found in the order of potency:
percent of dietary energy from non-staples (greatest impact)
access to sanitation and women's education
access to safe water
per capita dietary energy supply
Three of these determinants should receive attention in particular: access to sanitation, and diversity of calorie sources from food supplies. A study by the Institute of Development Studies has stressed that: "The first two should be prioritized because they have strong impacts yet are farthest below their desired levels".
The goal of UN agencies, governments and NGO is now to optimist nutrition during the first 1000 days of a child's life, from pregnancy to the child's second birthday, in order to reduce the prevalence of stunting. The first 1000 days in a child's life are a crucial "window of opportunity" because the brain develops rapidly, laying the foundation for future cognitive and social ability. Furthermore, it is also the time when young children are the most at risk of infections that lead to diarrhea. It is the time when they stop breast feeding (weaning process), begin to crawl, put things in their mouths and become exposed to fecal matter from open defecation and environmental enteropathies.
Dietary interventions to improve stunting
Previous interventions to reduce stunting have shown modest effects. Multiple micro-nutrient supplementation shows only small benefits for linear growth and results from studies supplementing lipid based nutrient supplements (LNS) to children are inconclusive. Educational interventions to improve complementary feeding may achieve behavioral change but have no or small effects on growth. Further, studies on the effect of micro-nutrient fortification, increased availability of key nutrients or increased energy density of complementary foods on stunting also show heterogeneous results. It is estimated that education interventions, if optimally designed and implemented, could reduce stunting by 0.6 z-scores while food-based interventions could reduce stunting by 0.5 z-scores, which is moderate compared to the average global growth deficit. Finally, the Lancet-series on maternal and child nutrition estimated that the impact of all existing interventions designed to improve nutrition and prevent related diseases in mothers and children, could reduce stunting at 3 years by merely 36%. Hence, factors explaining the shortfall in observed associations between child feeding practices and nutrient intake and linear growth, have increasingly been the focus of scientific interest.
Recent works showed promise that intervention with egg may improve linear growth in children. Comprehensive intervention package containing eggs also found to be effective in improving linear growth in children. However, the effect of egg intervention may not persist for longer period. Therefore, intervention programs should consider egg intervention for a longer period with emphasis on overall diet quality and improvement of environmental conditions.
Pregnant and lactating mothers
Ensuring proper nutrition of pregnant and lactating mothers is essential. Achieving so by helping women of reproductive age be in good nutritional status at conception is an excellent preventive measure. A focus on the pre-conception period has recently been introduced as a complement to the key phase of the 1000 days of pregnancy and first two years of life. An example of this are attempts to control anemia in women of reproductive age. A well-nourished mother is the first step of stunting prevention, decreasing chances of the baby being born of low birth-weight, which is the first risk factor for future malnutrition.
Balanced protein–energy supplementation in pregnancy seem to improve birth weight of children, with greater effects in undernourished women. Meanwhile, micronutrient supplements and lipid based nutrient supplements (LNS) (providing both macro-and micronutrients) during pregnancy have shown mixed effects on birth weight and -length. Similarly, studies supplementing LNS to mothers during pregnancy and lactation and their children during the complementary feeding period show heterogeneous results for stunting.
After birth, in terms of interventions for the child, early initiation of breastfeeding, together with exclusive breastfeeding for the first 6 months, are pillars of stunting prevention. Introducing proper complementary feeding after 6 months of age together with breastfeeding until age 2 is the next step.
Policy interventions
In summary, key policy interventions for the prevention of stunting are:
Improvement in nutrition surveillance activities to identify rates and trends of stunting and other forms of malnutrition within countries. This should be done with an equity perspective, as it is likely that stunting rates will vary greatly between different population groups. The most vulnerable should be prioritized. The same should be done for risk factors such as anemia, maternal under-nutrition, food insecurity, low birth-weight, breastfeeding practices etc. By collecting more detailed information, it is easier to ensure that policy interventions really address the root causes of stunting.
Political will to develop and implement national targets and strategies in line with evidence-based international guidelines as well as contextual factors.
Designing and implementing policies promoting nutritional and health well-being of mothers and women of reproductive age. The main focus should be on the 1000 days of pregnancy and first two years of life, but the pre-conception period should not be neglected as it can play a significant role in ensuring the fetus and baby's nutrition.
Designing and implementing policies promoting proper breastfeeding and complementary feeding practice (focusing on diet diversity for both macro and micronutrients). This can ensure optimal infant nutrition as well as protection from infections that can weaken the child's body. Labor policy ensuring mothers have the chance to breastfeed should be considered where necessary.
Introducing interventions addressing social and other health determinants of stunting, such as poor sanitation and access to drinking water, early marriages, intestinal parasite infections, malaria and other childhood preventable disease (referred to as “nutrition-sensitive interventions”), as well as the country's food security landscape. Interventions to keep adolescent girls in school can be effective at delaying marriage with subsequent nutritional benefits for both women and babies. Regulating milk substitutes is also very important to ensure that as many mothers as possible breastfeed their babies, unless a clear contraindication is present.
Broadly speaking, effective policies to reduce stunting require multisectoral approaches, strong political commitment, community involvement and integrated service delivery.
Epidemiology
According to the World Health organization if less than 20% of the population is affected by stunting, this is regarded as "low prevalence" in terms of public health significance. Values of 40% or more are regarded as very high prevalence, and values in between as medium to high prevalence.
UNICEF has estimated that: "Globally, more than one quarter (26 per cent) of children under 5 years of age were stunted in 2011 – roughly 165 million children worldwide." and "In sub-Saharan Africa, 40 per cent of children under 5 years of age are stunted; in South Asia, 39 per cent are stunted." The four countries with the highest prevalence are Timor-Leste, Burundi, Niger and Madagascar where more than half of children under 5 years old are stunted.
The 2020 edition of FAO's Near East and North Africa − Regional Overview of Food Security and Nutrition found that in 2019 22.5 percent of children under the age of five were stunted, 9.2 percent were wasted, and 9.9 percent were overweight across several Arab and North African countries.
Trends
As of 2015, it was estimated that there were 156 million stunted children under 5 in the world, 90% of them living in low and low-middle income countries. 56% of these were in Asia, and 37% in Africa. It is possible that some of these children concurrently had other forms of malnutrition, including wasting and stunting, and overweight and stunting. No statistics are currently available for these combined conditions. Stunting has been on the decline for the past 15 years. As a comparison, there were 255 million stunted children in 1990, 224 in 1995, 198 in 2000, 182 in 2005, 169 in 2010, and 156 in 2016. However, the decline of stunting is geographically uneven, and is unequal among different groups in society. A research paper published in January 2020, which mapped stunting, wasting and underweight in children in low- and middle-income countries, predicted that only five countries would meet global targets for reducing malnutrition by 2025 in all second administrative subdivisions.
Over the period 2000–2015, Asia reduced its stunting prevalence from 38 to 24%, Africa from 38 to 32%, and Latin America and the Caribbean from 18 to 11%. This equates to a relative reduction of 36, 17 and 39% respectively, indicating that Asia and Latin America and the Caribbean have displayed much larger improvements than Africa, which needs to address this issue with much more effort if it is to win the battle against a problem that has been crippling its development for decades. Of these regions, Latin America and the Caribbean are on track to achieve global targets set with global initiatives such as the United Nations Millennium Development Goals and the World Health Assembly targets (see following section on global targets).
Sub-regional stunting rates are as follows: In Africa, the highest rates are observed in East Africa (37.5%). All other Sub-Saharan sub-regions also have high rates, with 32.1% in West Africa, 31.2% in Central Africa, and 28.4% in Southern Africa. North Africa is at 18%, and the Middle East at 16.2%. In Asia, the highest rate is observed in South Asia at 34.4%. South-East Asia is at 26.3%. Pacific Islands also display a high rate at 38.2%. Central and South America are respectively at 15.6 and 9.9%. South Asia, given its very high population at over 1 billion and high prevalence rate of stunting, is the region currently hosting the highest absolute number of children with stunting.
The sheer number of stunted children has increased in Africa from 50.4 to 58.5 million in the time 2000–2015. This is despite the reduction in percentage prevalence of stunting, and is due to the high rates of population growth. The data therefore indicates that the rate of reduction of stunting in Africa has not been able to counterbalance the increased number of growing children that fall into the trap of malnutrition, due to population growth in the region, creating a cycle. This is also true in Oceania, unlike Asia and Latin America and the Caribbean where substantial absolute reductions in the number of stunted children have been observed (for example, Asia reduced its number of stunted children from 133 million to 88 million between 2000 and 2015).
The reduction in stunting is closely linked to poverty reduction and the will and ability of governments to set up solid multisectoral approaches to reduce chronic malnutrition. Low income countries are the only group with more stunted children today than in the year 2000. Conversely, all other countries (high-income, upper-middle income, lower-middle income) have achieved reductions in the numbers of stunted children. This perpetuates a vicious cycle of poverty and malnutrition, whereby malnourished children are not able to maximally contribute to economic development as adults, and poverty increases chances of malnutrition.
Research
The Water and Sanitation Program of the World Bank has investigated links between lack of sanitation and stunting in Vietnam and Lao PDR. An example is in Vietnam where the lack of sanitation in rural villages in mountainous regions of Vietnam led to five-year-old children being 3.7 cm shorter than healthy children living in villages with good access to sanitation. This difference in height is irreversible and matters a great deal for a child's cognitive development and future productive potential.
Review articles
The Lancet has published two comprehensive series on maternal and child nutrition, in 2008 and 2013. The series review the epidemiology of global malnutrition and analyze the state of the evidence for cost-effective interventions that should be scaled-up to achieve impact and global targets. In the first of such series, investigators define the importance of the 1000 day and identify child malnutrition as being responsible for one third of all child deaths worldwide. This finding is key in that it points at malnutrition as a key determinant of child mortality that is often overlooked. When a child dies of pneumonia, malaria or diarrhea (some of the causes of child mortality in the world), it may well be that malnutrition is a key contributing factor that prevents the body from successfully fighting the infection and recovering from these diseases.
In the follow-up series in 2013, the focus on undernutrition is expanded to the increasing burden of obesity in both high, middle and low income countries. Several countries with high levels of child stunting and undernutrition are starting to display worrisome increasing trends of child obesity concurrently, due to increased wealth and the persistence of significant inequalities. The challenges these countries face are particularly difficult as they require intervening on two levels on what has come to be called “double burden of malnutrition”. As an example, in India 30% of children under 5 years of age are stunted, and 20% are overweight. Neglecting these nutritional problems is not an option anymore if countries are to escape poverty traps and provide opportunities to their people to live fulfilling productive lives without stunting.
Nutritional interventions such as dietary supplementation and nutritional education have the potential to decrease stunting.
Examples
The 2012 World Health Assembly, with its 194 member states, convened to discuss global issues of maternal, infant and young child nutrition, and developed a plan with 6 targets for 2025. The first of such targets aims to reduce by 40% the number of children who are stunted in the world, by 2025. This would correspond to 100 million stunted children in 2025. At the current reduction rate, the predicted number in 2025 will be 127 million, indicating the need to scale-up and intensify efforts if the global community is to reach its goals.
The World Bank estimates that the extra cost to achieve the reduction goal will be $8.50 yearly per stunted child, for a total of $49.6 Billion for the next decade. Stunting has been shown to be one of the most cost-effective global health problems to invest in, with an estimated return on investment of $18 for every dollar spent thanks to its impact on economic productivity. Despite the evidence in favor of investing in the reduction of stunting, current investments are too low at about $2.9 billion per year, with $1.6 billion coming from Governments, $0.2 billion from donors, and $1.1 paid by individuals.
Sustainable Development Goals
In 2015, the United Nations and its member states agreed on a new sustainable development agenda to promote prosperity and reduce poverty, putting forward 17 Sustainable Development Goals (SDGs) to be achieved by 2030. SDG 2 aims to “End hunger, achieve food security and improved nutrition, and promote sustainable agriculture”. Sub-goal 2.2. aims to “by 2030 end all forms of malnutrition, including achieving by 2025 the internationally agreed targets on stunting and wasting in children under five years of age, and address the nutritional needs of adolescent girls, pregnant and lactating women, and older persons”.
The global community has recognized more and more the critical importance of stunting during the past decade. Investments to address it have increased but remain far from being sufficient to solve it and unleash the human potential that remains trapped in malnutrition.
The "Scaling Up Nutrition Movement (SUN)" movement is the main network of governments, non-governmental and international organizations, donors, private companies and academic institutions working together in pursuit of improved global nutrition and a world without hunger and malnutrition. It was launched at the UN General Assembly of 2010 and it calls for country-led multi-sectoral strategies to address child malnutrition by scaling-up evidence-based interventions in both nutrition specific and sensitive areas. As of 2016, 50 countries have joined the SUN Movement with strategies that are aligned with international frameworks of action.
Brazil
Brazil displayed a remarkable reduction in the rates of child stunting under age 5, from 37% in 1974, to 7.1% in 2007. This happened in association with impressive social and economic development that reduced the numbers of Brazilians living in extreme poverty (less than $1.25 per day) from 25.6% in 1990 to 4.8% in 2008. The successful reduction in child malnutrition in Brazil can be attributed to strong political commitment that led to improvements in the water and sanitation system, increased female schooling, scale-up of quality maternal and child health services, increased economic power at family level (including successful cash transfer programs), and improvements in food security throughout the country.
Bangladesh
Nearly one-third of the children under five years of age are stunted in Bangladesh and 9% are severely stunted. The country is on track in reducing the prevalence of stunted growth. If the current trend continues, the prevalence would be 21% in 2025, while the target is 27%. Maternal undernutrition and increased pathogen load in the intestine are the major risk factors of stunting in Bangladeshi children. Daily supplementation with egg, cow milk, and micronutrient powder found to be effective in improving linear growth of children in a community-based trial in Bangladesh.
Peru
After a decade (1995–2005) in which stunting rates stagnated in the country, Peru designed and implemented a national strategy against child malnutrition called crecer ("grow"), which complemented a social development conditional cash-transfer program called juntos, which included a nutritional component. The strategy was multisectoral in that it involved the health, education, water, sanitation and hygiene, agriculture and housing sectors and stakeholders. It was led by the Government and the Prime Minister himself, and included non-governmental partners at both central, regional and community level. After the strategy was implemented, stunting went from 22.9% to 17.9% (2005–2010), with very significant improvements in rural areas where it had been more difficult to reduce stunting rates in the past.
India (Maharashtra)
The State of Maharashtra in Central-Western India has been able to produce an impressive reduction in stunting rates in children under 2 years of age from 44% to 22.8% in the 2005–2012 period. This is particularly remarkable given the immense challenges India has faced to address malnutrition, and that the country hosts almost half of all stunted children under 5 in the world. This was achieved through integrated community-based programs that were designed by a central advisory body that promoted multisectoral collaboration, provided advice to policy-makers on evidence-based solutions, and advocated for the key role of the 1000 days (pregnancy and first two years of life).
Nepal
In Nepal, short maternal stature, low maternal education, poor access to health services and poverty are strong determinants for stunting. However, in Nepal, stunting has decreased from 57% in 2001 to 36% in 2016, with lower prevalence in urban than in rural settings.
Philippines
In the Philippines, one in three children below five years old is stunted. Even though the country's economic growth has steadily increased by 4% annually, almost a third of Filipino children have stunted growth. The prevalence of stunting declined during the early 2000s but has remained the same since then, with the 2019 rate (28.8 percent) only marginally lower than that of 2008. Researchers attribute the problem to micro-nutrient deficiencies brought on by poverty, maternal under-education, food insecurity, and poor environmental conditions. To address stunting and other health and food security issues, the Philippine Plan of Action for Nutrition (PPAN) was established as an umbrella initiative to meet health and nutrition targets in the country by 2028. Since 2015, there has been a decline in stunting across all age groups, from infants to teenagers, with the most significant improvement observed among 5 to 10-year-olds, dropping from 31.2 percent in 2015 to 19.7 percent in 2021.
| Biology and health sciences | Health and fitness: General | Health |
3489162 | https://en.wikipedia.org/wiki/Oriental%20garden%20lizard | Oriental garden lizard | The oriental garden lizard (Calotes versicolor), also called the eastern garden lizard, Indian garden lizard, common garden lizard, bloodsucker or changeable lizard, is an agamid lizard found widely distributed in Indo-Malaya. It has also been introduced in many other parts of the world.
Description
Calotes versicolor is an insectivore, and the male gets a bright red throat in the breeding season.
It measures over 10 cm (3.9 in) in length snout-to-vent. Total length including the tail is up to 37 cm (14.5 in).
Two small groups of spines, perfectly separated from each other, above each tympanum. Dorsal crest moderately elevated on the neck and anterior part of the trunk, extending on to the root of the tail in large individuals, and gradually disappearing on the middle of the trunk in younger ones. No fold in front of the shoulder, but the scales behind the lower jaw are much smaller than the others; gular sac not developed. From thirty-nine to forty-three series of scales round the middle of the trunk. The hind foot (measured from the heel to the extremity of the fourth toe) is not much longer than the head in the adult, whilst it is considerably longer in the young. The coloration is very variable, sometimes uniform brownish or greyish-olive or yellowish. Generally broad brown bands across the back, interrupted by a yellowish lateral band. Black streaks radiate from the eye, and some of them are continued over the throat, running obliquely backwards, belly frequently with greyish longitudinal stripes, one along the median line being the most distinct; young and half-grown specimens have a dark, black-edged band across the inter-orbital region.
The ground-colour is generally a light brownish olive, but the lizard can change it to bright red, to black, and to a mixture of both. This change is sometimes confined to the head, at other times diffused over the whole body and tail. A common state in which it may be seen (as stated by Mr. Jerdon) is, seated on a hedge or bush, with the tail and limbs black, head and neck yellow picked out with red, and the rest of the body red. Jerdon and Blyth agree that these bright, changeable colours are peculiar to the male during the breeding-season, which falls in the months of May to early October.
C. A. L. Guenther mentioned that Mouhot had collected in Siam one of those fine variations of colours, which, however, appear to be infinite. It has the usual cross streaks between the eyes and the radiating lines continent of India to China; it is very common in Ceylon, not extending into the temperate zone of the Himalayas. Ceylonese specimens are generally somewhat larger; one of them measured 16 inches, the tail taking 11 inches. It is found in hedges and trees; it is known in Ceylon under the name of "Bloodsucker", a designation the origin of which cannot be satisfactorily traced; in the opinion of Kelaart, the name was given to it from the occasional reddish hue of the throat and neck. "Roktochosha (রক্তচোষা)" is also a local name in the Bengali language, which also translates to "bloodsucker".
The female lays from five to sixteen soft oval eggs, about 5/8 of an inch long, in hollows of trees, or in holes in the soil which they have burrowed, afterward covering them up. The young appear in about eight or nine weeks. In a hot sunny day a solitary bloodsucker may be seen on a twig or on a wall, basking in the sun, with mouth wide open. After a shower of rain numbers of them are seen to come down on the ground and pick up the larva and small insects which fall from the trees during the showers. Changeable lizards escape danger by darting to the nearest tree. If the predator comes even closer, they will scale to the side of the tree facing away from the predator and very swiftly dart up the tree. The predator looks behind the tree only to see that the lizard is up in the branches.
During the breeding season, the male's head and shoulders turns bright orange to crimson and his throat black. Males also turn red-headed after a successful battle with rivals. Both males and females have a crest from the head to nearly the tail, hence their other common name, "crested tree lizard".
Unlike some other lizards, they do not drop their tails (autotomy), and their tails can be very long, stiff and pointy. Like other reptiles, they shed their skins. Like chameleons, changeable lizards can move each of their eyes in different directions.
Distribution
The native range of the species includes southeastern Iran, Afghanistan, Bangladesh, Bhutan, Cambodia, India (including the Andaman Islands), Indonesia (Sumatra), Malaysia (western), Maldives, Mauritius (Reunion, Rodrigues), Myanmar, Nepal, Pakistan, Philippines, Sri Lanka (Ceylon), Thailand, Vietnam (including Pulo Condore Island). It has been introduced to Brunei, Celebes, Oman, Seychelles, Singapore and United States. The lizards were introduced to Singapore from Malaysia and Thailand in the 1980s. In Singapore, they are a threat to the native green-crested lizard.
The changeable lizard is relatively common and found in a wide range of habitats. They appear to adapt well to humans and are thus not endangered. They are commonly found among undergrowth, in open habitats as well as highly urban areas. However, in China people regularly kill them, as they are viewed as pests.
Diet
Changeable lizards eat mainly insects such as crickets, grasshoppers, and ants; as well as small vertebrates, including rodents and other lizards including common house geckos and day geckos. They have teeth which are designed for gripping prey and not tearing, and thus they usually shake prey to stun it then swallow it whole. Sometimes, young, inexperienced changeable lizards may choke on prey that is too large. Changeable lizards also occasionally consume vegetable matter.
Reproduction
Males become highly territorial during breeding season. They discourage intruding males by brightening their red heads and doing "push-ups". Each tries to attract a female by inflating his throat and drawing attention to his handsomely colored head. Oviparous; about 10—20 eggs are laid, buried in moist soil. The eggs are long, spindle-shaped and covered with a leathery skin. They hatch in about 6–7 weeks. They are able to breed at about 1 year old.
| Biology and health sciences | Iguania | Animals |
113079 | https://en.wikipedia.org/wiki/Winnowing | Winnowing | Winnowing is a process by which chaff is separated from grain. It can also be used to remove pests from stored grain. Winnowing usually follows threshing in grain preparation. In its simplest form, it involves throwing the mixture into the air so that the wind blows away the lighter chaff, while the heavier grains fall back down for recovery. Techniques included using a winnowing fan (a shaped basket shaken to raise the chaff) or using a tool (a winnowing fork or shovel) on a pile of harvested grain.
In Greek culture
The winnowing-fan (λίκνον [líknon], also meaning a "cradle") featured in the rites accorded Dionysus and in the Eleusinian Mysteries: "it was a simple agricultural implement taken over and mysticized by the religion of Dionysus," Jane Ellen Harrison remarked. Dionysus Liknites ("Dionysus of the winnowing fan") was wakened by the Dionysian women, in this instance called Thyiades, in a cave on Parnassus high above Delphi; the winnowing-fan links the god connected with the mystery religions to the agricultural cycle, but mortal Greek babies too were laid in a winnowing-fan. In Callimachus's Hymn to Zeus, Adrasteia lays the infant Zeus in a golden líknon, her goat suckles him and he is given honey. In the Odyssey, the dead oracle Teiresias tells Odysseus to walk away from Ithaca with an oar until a wayfarer tells him it is a winnowing fan (i.e., until Odysseus has come so far from the sea that people don't recognize oars), and there to build a shrine to Poseidon.
China
In ancient China, the method was improved by mechanization with the development of the rotary winnowing fan, which used a cranked fan to produce the airstream. This was featured in Wang Zhen's book the Nong Shu of 1313 AD.
In Europe
In Saxon settlements such as one identified in Northumberland as Bede's Ad Gefrin (now called Yeavering) the buildings were shown by an excavator's reconstruction to have opposed entries. In barns a draught created by the use of these opposed doorways was used in winnowing.
The technique developed by the Chinese was not adopted in Europe until the 18th century when winnowing machines used a 'sail fan'. The rotary winnowing fan was exported to Europe, brought there by Dutch sailors between 1700 and 1720. Apparently, they had obtained them from the Dutch settlement of Batavia in Java, Dutch East Indies. The Swedes imported some from south China at about the same time and Jesuits had taken several to France from China by 1720. Until the beginning of the 18th century, no rotary winnowing fans existed in the West.
In the United States
The development of the winnowing barn allowed rice plantations in South Carolina to increase their yields dramatically.
Mechanization of the process
In 1737 Andrew Rodger, a farmer on the estate of Cavers in Roxburghshire, developed a winnowing machine for corn, called a 'Fanner'. These were successful and the family sold them throughout Scotland for many years. Some Scottish Presbyterian ministers saw the fanners as sins against God, for the wind was a thing specially made by him and an artificial wind was a daring and impious attempt to usurp what belonged to God alone. As the Industrial Revolution progressed, the winnowing process was mechanized by the invention of additional winnowing machines, such as fanning mills.
| Technology | Agronomical techniques | null |
113087 | https://en.wikipedia.org/wiki/System%20of%20linear%20equations | System of linear equations | In mathematics, a system of linear equations (or linear system) is a collection of two or more linear equations involving the same variables.
For example,
is a system of three equations in the three variables . A solution to a linear system is an assignment of values to the variables such that all the equations are simultaneously satisfied. In the example above, a solution is given by the ordered triple
since it makes all three equations valid.
Linear systems are a fundamental part of linear algebra, a subject used in most modern mathematics. Computational algorithms for finding the solutions are an important part of numerical linear algebra, and play a prominent role in engineering, physics, chemistry, computer science, and economics. A system of non-linear equations can often be approximated by a linear system (see linearization), a helpful technique when making a mathematical model or computer simulation of a relatively complex system.
Very often, and in this article, the coefficients and solutions of the equations are constrained to be real or complex numbers, but the theory and algorithms apply to coefficients and solutions in any field. For other algebraic structures, other theories have been developed. For coefficients and solutions in an integral domain, such as the ring of integers, see Linear equation over a ring. For coefficients and solutions that are polynomials, see Gröbner basis. For finding the "best" integer solutions among many, see Integer linear programming. For an example of a more exotic structure to which linear algebra can be applied, see Tropical geometry.
Elementary examples
Trivial example
The system of one equation in one unknown
has the solution
However, most interesting linear systems have at least two equations.
Simple nontrivial example
The simplest kind of nontrivial linear system involves two equations and two variables:
One method for solving such a system is as follows. First, solve the top equation for in terms of :
Now substitute this expression for x into the bottom equation:
This results in a single equation involving only the variable . Solving gives , and substituting this back into the equation for yields . This method generalizes to systems with additional variables (see "elimination of variables" below, or the article on elementary algebra.)
General form
A general system of m linear equations with n unknowns and coefficients can be written as
where are the unknowns, are the coefficients of the system, and are the constant terms.
Often the coefficients and unknowns are real or complex numbers, but integers and rational numbers are also seen, as are polynomials and elements of an abstract algebraic structure.
Vector equation
One extremely helpful view is that each unknown is a weight for a column vector in a linear combination.
This allows all the language and theory of vector spaces (or more generally, modules) to be brought to bear. For example, the collection of all possible linear combinations of the vectors on the left-hand side is called their span, and the equations have a solution just when the right-hand vector is within that span. If every vector within that span has exactly one expression as a linear combination of the given left-hand vectors, then any solution is unique. In any event, the span has a basis of linearly independent vectors that do guarantee exactly one expression; and the number of vectors in that basis (its dimension) cannot be larger than m or n, but it can be smaller. This is important because if we have m independent vectors a solution is guaranteed regardless of the right-hand side, and otherwise not guaranteed.
Matrix equation
The vector equation is equivalent to a matrix equation of the form
where A is an m×n matrix, x is a column vector with n entries, and b is a column vector with m entries.
The number of vectors in a basis for the span is now expressed as the rank of the matrix.
Solution set
A solution of a linear system is an assignment of values to the variables
such that each of the equations is satisfied. The set of all possible solutions is called the solution set.
A linear system may behave in any one of three possible ways:
The system has infinitely many solutions.
The system has a unique solution.
The system has no solution.
Geometric interpretation
For a system involving two variables (x and y), each linear equation determines a line on the xy-plane. Because a solution to a linear system must satisfy all of the equations, the solution set is the intersection of these lines, and is hence either a line, a single point, or the empty set.
For three variables, each linear equation determines a plane in three-dimensional space, and the solution set is the intersection of these planes. Thus the solution set may be a plane, a line, a single point, or the empty set. For example, as three parallel planes do not have a common point, the solution set of their equations is empty; the solution set of the equations of three planes intersecting at a point is single point; if three planes pass through two points, their equations have at least two common solutions; in fact the solution set is infinite and consists in all the line passing through these points.
For n variables, each linear equation determines a hyperplane in n-dimensional space. The solution set is the intersection of these hyperplanes, and is a flat, which may have any dimension lower than n.
General behavior
In general, the behavior of a linear system is determined by the relationship between the number of equations and the number of unknowns. Here, "in general" means that a different behavior may occur for specific values of the coefficients of the equations.
In general, a system with fewer equations than unknowns has infinitely many solutions, but it may have no solution. Such a system is known as an underdetermined system.
In general, a system with the same number of equations and unknowns has a single unique solution.
In general, a system with more equations than unknowns has no solution. Such a system is also known as an overdetermined system.
In the first case, the dimension of the solution set is, in general, equal to , where n is the number of variables and m is the number of equations.
The following pictures illustrate this trichotomy in the case of two variables:
{| border=0 cellpadding=5
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| width="150" align="center" |
| width="150" align="center" |
|-
| align="center" | One equation
| align="center" | Two equations
| align="center" | Three equations
|}
The first system has infinitely many solutions, namely all of the points on the blue line. The second system has a single unique solution, namely the intersection of the two lines. The third system has no solutions, since the three lines share no common point.
It must be kept in mind that the pictures above show only the most common case (the general case). It is possible for a system of two equations and two unknowns to have no solution (if the two lines are parallel), or for a system of three equations and two unknowns to be solvable (if the three lines intersect at a single point).
A system of linear equations behave differently from the general case if the equations are linearly dependent, or if it is inconsistent and has no more equations than unknowns.
Properties
Independence
The equations of a linear system are independent if none of the equations can be derived algebraically from the others. When the equations are independent, each equation contains new information about the variables, and removing any of the equations increases the size of the solution set. For linear equations, logical independence is the same as linear independence.
For example, the equations
are not independent — they are the same equation when scaled by a factor of two, and they would produce identical graphs. This is an example of equivalence in a system of linear equations.
For a more complicated example, the equations
are not independent, because the third equation is the sum of the other two. Indeed, any one of these equations can be derived from the other two, and any one of the equations can be removed without affecting the solution set. The graphs of these equations are three lines that intersect at a single point.
Consistency
A linear system is inconsistent if it has no solution, and otherwise, it is said to be consistent. When the system is inconsistent, it is possible to derive a contradiction from the equations, that may always be rewritten as the statement .
For example, the equations
are inconsistent. In fact, by subtracting the first equation from the second one and multiplying both sides of the result by 1/6, we get . The graphs of these equations on the xy-plane are a pair of parallel lines.
It is possible for three linear equations to be inconsistent, even though any two of them are consistent together. For example, the equations
are inconsistent. Adding the first two equations together gives , which can be subtracted from the third equation to yield . Any two of these equations have a common solution. The same phenomenon can occur for any number of equations.
In general, inconsistencies occur if the left-hand sides of the equations in a system are linearly dependent, and the constant terms do not satisfy the dependence relation. A system of equations whose left-hand sides are linearly independent is always consistent.
Putting it another way, according to the Rouché–Capelli theorem, any system of equations (overdetermined or otherwise) is inconsistent if the rank of the augmented matrix is greater than the rank of the coefficient matrix. If, on the other hand, the ranks of these two matrices are equal, the system must have at least one solution. The solution is unique if and only if the rank equals the number of variables. Otherwise the general solution has k free parameters where k is the difference between the number of variables and the rank; hence in such a case there is an infinitude of solutions. The rank of a system of equations (that is, the rank of the augmented matrix) can never be higher than [the number of variables] + 1, which means that a system with any number of equations can always be reduced to a system that has a number of independent equations that is at most equal to [the number of variables] + 1.
Equivalence
Two linear systems using the same set of variables are equivalent if each of the equations in the second system can be derived algebraically from the equations in the first system, and vice versa. Two systems are equivalent if either both are inconsistent or each equation of each of them is a linear combination of the equations of the other one. It follows that two linear systems are equivalent if and only if they have the same solution set.
Solving a linear system
There are several algorithms for solving a system of linear equations.
Describing the solution
When the solution set is finite, it is reduced to a single element. In this case, the unique solution is described by a sequence of equations whose left-hand sides are the names of the unknowns and right-hand sides are the corresponding values, for example . When an order on the unknowns has been fixed, for example the alphabetical order the solution may be described as a vector of values, like for the previous example.
To describe a set with an infinite number of solutions, typically some of the variables are designated as free (or independent, or as parameters), meaning that they are allowed to take any value, while the remaining variables are dependent on the values of the free variables.
For example, consider the following system:
The solution set to this system can be described by the following equations:
Here z is the free variable, while x and y are dependent on z. Any point in the solution set can be obtained by first choosing a value for z, and then computing the corresponding values for x and y.
Each free variable gives the solution space one degree of freedom, the number of which is equal to the dimension of the solution set. For example, the solution set for the above equation is a line, since a point in the solution set can be chosen by specifying the value of the parameter z. An infinite solution of higher order may describe a plane, or higher-dimensional set.
Different choices for the free variables may lead to different descriptions of the same solution set. For example, the solution to the above equations can alternatively be described as follows:
Here x is the free variable, and y and z are dependent.
Elimination of variables
The simplest method for solving a system of linear equations is to repeatedly eliminate variables. This method can be described as follows:
In the first equation, solve for one of the variables in terms of the others.
Substitute this expression into the remaining equations. This yields a system of equations with one fewer equation and unknown.
Repeat steps 1 and 2 until the system is reduced to a single linear equation.
Solve this equation, and then back-substitute until the entire solution is found.
For example, consider the following system:
Solving the first equation for x gives , and plugging this into the second and third equation yields
Since the LHS of both of these equations equal y, equating the RHS of the equations. We now have:
Substituting z = 2 into the second or third equation gives y = 8, and the values of y and z into the first equation yields x = −15. Therefore, the solution set is the ordered triple .
Row reduction
In row reduction (also known as Gaussian elimination), the linear system is represented as an augmented matrix
This matrix is then modified using elementary row operations until it reaches reduced row echelon form. There are three types of elementary row operations:
Type 1: Swap the positions of two rows.
Type 2: Multiply a row by a nonzero scalar.
Type 3: Add to one row a scalar multiple of another.
Because these operations are reversible, the augmented matrix produced always represents a linear system that is equivalent to the original.
There are several specific algorithms to row-reduce an augmented matrix, the simplest of which are Gaussian elimination and Gauss–Jordan elimination. The following computation shows Gauss–Jordan elimination applied to the matrix above:
The last matrix is in reduced row echelon form, and represents the system , , . A comparison with the example in the previous section on the algebraic elimination of variables shows that these two methods are in fact the same; the difference lies in how the computations are written down.
Cramer's rule
Cramer's rule is an explicit formula for the solution of a system of linear equations, with each variable given by a quotient of two determinants. For example, the solution to the system
is given by
For each variable, the denominator is the determinant of the matrix of coefficients, while the numerator is the determinant of a matrix in which one column has been replaced by the vector of constant terms.
Though Cramer's rule is important theoretically, it has little practical value for large matrices, since the computation of large determinants is somewhat cumbersome. (Indeed, large determinants are most easily computed using row reduction.)
Further, Cramer's rule has very poor numerical properties, making it unsuitable for solving even small systems reliably, unless the operations are performed in rational arithmetic with unbounded precision.
Matrix solution
If the equation system is expressed in the matrix form , the entire solution set can also be expressed in matrix form. If the matrix A is square (has m rows and n=m columns) and has full rank (all m rows are independent), then the system has a unique solution given by
where is the inverse of A. More generally, regardless of whether m=n or not and regardless of the rank of A, all solutions (if any exist) are given using the Moore–Penrose inverse of A, denoted , as follows:
where is a vector of free parameters that ranges over all possible n×1 vectors. A necessary and sufficient condition for any solution(s) to exist is that the potential solution obtained using satisfy — that is, that If this condition does not hold, the equation system is inconsistent and has no solution. If the condition holds, the system is consistent and at least one solution exists. For example, in the above-mentioned case in which A is square and of full rank, simply equals and the general solution equation simplifies to
as previously stated, where has completely dropped out of the solution, leaving only a single solution. In other cases, though, remains and hence an infinitude of potential values of the free parameter vector give an infinitude of solutions of the equation.
Other methods
While systems of three or four equations can be readily solved by hand (see Cracovian), computers are often used for larger systems. The standard algorithm for solving a system of linear equations is based on Gaussian elimination with some modifications. Firstly, it is essential to avoid division by small numbers, which may lead to inaccurate results. This can be done by reordering the equations if necessary, a process known as pivoting. Secondly, the algorithm does not exactly do Gaussian elimination, but it computes the LU decomposition of the matrix A. This is mostly an organizational tool, but it is much quicker if one has to solve several systems with the same matrix A but different vectors b.
If the matrix A has some special structure, this can be exploited to obtain faster or more accurate algorithms. For instance, systems with a symmetric positive definite matrix can be solved twice as fast with the Cholesky decomposition. Levinson recursion is a fast method for Toeplitz matrices. Special methods exist also for matrices with many zero elements (so-called sparse matrices), which appear often in applications.
A completely different approach is often taken for very large systems, which would otherwise take too much time or memory. The idea is to start with an initial approximation to the solution (which does not have to be accurate at all), and to change this approximation in several steps to bring it closer to the true solution. Once the approximation is sufficiently accurate, this is taken to be the solution to the system. This leads to the class of iterative methods. For some sparse matrices, the introduction of randomness improves the speed of the iterative methods. One example of an iterative method is the Jacobi method, where the matrix is split into its diagonal component and its non-diagonal component . An initial guess is used at the start of the algorithm. Each subsequent guess is computed using the iterative equation:
When the difference between guesses and is sufficiently small, the algorithm is said to have converged on the solution.
There is also a quantum algorithm for linear systems of equations.
Homogeneous systems
A system of linear equations is homogeneous if all of the constant terms are zero:
A homogeneous system is equivalent to a matrix equation of the form
where A is an matrix, x is a column vector with n entries, and 0 is the zero vector with m entries.
Homogeneous solution set
Every homogeneous system has at least one solution, known as the zero (or trivial) solution, which is obtained by assigning the value of zero to each of the variables. If the system has a non-singular matrix () then it is also the only solution. If the system has a singular matrix then there is a solution set with an infinite number of solutions. This solution set has the following additional properties:
If u and v are two vectors representing solutions to a homogeneous system, then the vector sum is also a solution to the system.
If u is a vector representing a solution to a homogeneous system, and r is any scalar, then ru is also a solution to the system.
These are exactly the properties required for the solution set to be a linear subspace of Rn. In particular, the solution set to a homogeneous system is the same as the null space of the corresponding matrix A.
Relation to nonhomogeneous systems
There is a close relationship between the solutions to a linear system and the solutions to the corresponding homogeneous system:
Specifically, if p is any specific solution to the linear system , then the entire solution set can be described as
Geometrically, this says that the solution set for is a translation of the solution set for . Specifically, the flat for the first system can be obtained by translating the linear subspace for the homogeneous system by the vector p.
This reasoning only applies if the system has at least one solution. This occurs if and only if the vector b lies in the image of the linear transformation A.
| Mathematics | Linear algebra | null |
113222 | https://en.wikipedia.org/wiki/Law%20of%20sines | Law of sines | In trigonometry, the law of sines, sine law, sine formula, or sine rule is an equation relating the lengths of the sides of any triangle to the sines of its angles. According to the law,
where , and are the lengths of the sides of a triangle, and , and are the opposite angles (see figure 2), while is the radius of the triangle's circumcircle. When the last part of the equation is not used, the law is sometimes stated using the reciprocals;
The law of sines can be used to compute the remaining sides of a triangle when two angles and a side are known—a technique known as triangulation. It can also be used when two sides and one of the non-enclosed angles are known. In some such cases, the triangle is not uniquely determined by this data (called the ambiguous case) and the technique gives two possible values for the enclosed angle.
The law of sines is one of two trigonometric equations commonly applied to find lengths and angles in scalene triangles, with the other being the law of cosines.
The law of sines can be generalized to higher dimensions on surfaces with constant curvature.
History
The equivalent of the law of sines, that the sides of a triangle are proportional to the chords of double the opposite angles, was known to the 2nd century Hellenistic astronomer Ptolemy and used occasionally in his Almagest.
Statements related to the law of sines appear in the astronomical and trigonometric work of 7th century Indian mathematician Brahmagupta. In his Brāhmasphuṭasiddhānta, Brahmagupta expresses the circumradius of a triangle as the product of two sides divided by twice the altitude; the law of sines can be derived by alternately expressing the altitude as the sine of one or the other base angle times its opposite side, then equating the two resulting variants. An equation even closer to the modern law of sines appears in Brahmagupta's Khaṇḍakhādyaka, in a method for finding the distance between the Earth and a planet following an epicycle; however, Brahmagupta never treated the law of sines as an independent subject or used it more systematically for solving triangles.
The spherical law of sines is sometimes credited to 10th century scholars Abu-Mahmud Khujandi or Abū al-Wafāʾ (it appears in his Almagest), but it is given prominence in Abū Naṣr Manṣūr's Treatise on the Determination of Spherical Arcs, and was credited to Abū Naṣr Manṣūr by his student al-Bīrūnī in his Keys to Astronomy. Ibn Muʿādh al-Jayyānī's 11th-century Book of Unknown Arcs of a Sphere also contains the spherical law of sines.
The 13th-century Persian mathematician Naṣīr al-Dīn al-Ṭūsī stated and proved the planar law of sines: In any plane triangle, the ratio of the sides is equal to the ratio of the sines of the angles opposite to those sides. That is, in triangle ABC, we have AB : AC = Sin(∠ACB) : Sin(∠ABC) By employing the law of sines, al-Tusi could solve triangles where either two angles and a side were known or two sides and an angle opposite one of them were given. For triangles with two sides and the included angle, he divided them into right triangles that he could then solve. When three sides were given, he dropped a perpendicular line and then used Proposition II-13 of Euclid's Elements (a geometric version of the law of cosines). Al-Tusi established the important result that if the sum or difference of two arcs is provided along with the ratio of their sines, then the arcs can be calculated.
According to Glen Van Brummelen, "The Law of Sines is really Regiomontanus's foundation for his solutions of right-angled triangles in Book IV, and these solutions are in turn the bases for his solutions of general triangles." Regiomontanus was a 15th-century German mathematician.
Proof
With the side of length as the base, the triangle's altitude can be computed as or as . Equating these two expressions gives
and similar equations arise by choosing the side of length or the side of length as the base of the triangle.
The ambiguous case of triangle solution
When using the law of sines to find a side of a triangle, an ambiguous case occurs when two separate triangles can be constructed from the data provided (i.e., there are two different possible solutions to the triangle). In the case shown below they are triangles and .
Given a general triangle, the following conditions would need to be fulfilled for the case to be ambiguous:
The only information known about the triangle is the angle and the sides and .
The angle is acute (i.e., < 90°).
The side is shorter than the side (i.e., ).
The side is longer than the altitude from angle , where (i.e., ).
If all the above conditions are true, then each of angles and produces a valid triangle, meaning that both of the following are true:
From there we can find the corresponding and or and if required, where is the side bounded by vertices and and is bounded by and .
Examples
The following are examples of how to solve a problem using the law of sines.
Example 1
Given: side , side , and angle . Angle is desired.
Using the law of sines, we conclude that
Note that the potential solution is excluded because that would necessarily give .
Example 2
If the lengths of two sides of the triangle and are equal to , the third side has length , and the angles opposite the sides of lengths , , and are , , and respectively then
Relation to the circumcircle
In the identity
the common value of the three fractions is actually the diameter of the triangle's circumcircle. This result dates back to Ptolemy.
Proof
As shown in the figure, let there be a circle with inscribed and another inscribed that passes through the circle's center . The has a central angle of and thus by Thales's theorem. Since is a right triangle,
where is the radius of the circumscribing circle of the triangle. Angles and lie on the same circle and subtend the same chord ; thus, by the inscribed angle theorem, Therefore,
Rearranging yields
Repeating the process of creating with other points gives
Relationship to the area of the triangle
The area of a triangle is given by where is the angle enclosed by the sides of lengths and . Substituting the sine law into this equation gives
Taking as the circumscribing radius,
It can also be shown that this equality implies
where is the area of the triangle and is the semiperimeter
The second equality above readily simplifies to Heron's formula for the area.
The sine rule can also be used in deriving the following formula for the triangle's area: denoting the semi-sum of the angles' sines as we have
where is the radius of the circumcircle:
Spherical law of sines
The spherical law of sines deals with triangles on a sphere, whose sides are arcs of great circles.
Suppose the radius of the sphere is 1. Let , , and be the lengths of the great-arcs that are the sides of the triangle. Because it is a unit sphere, , , and are the angles at the center of the sphere subtended by those arcs, in radians. Let , , and be the angles opposite those respective sides. These are dihedral angles between the planes of the three great circles.
Then the spherical law of sines says:
Vector proof
Consider a unit sphere with three unit vectors , and drawn from the origin to the vertices of the triangle. Thus the angles , , and are the angles , , and , respectively. The arc subtends an angle of magnitude at the centre. Introduce a Cartesian basis with along the -axis and in the -plane making an angle with the -axis. The vector projects to in the -plane and the angle between and the -axis is . Therefore, the three vectors have components:
The scalar triple product, is the volume of the parallelepiped formed by the position vectors of the vertices of the spherical triangle , and . This volume is invariant to the specific coordinate system used to represent , and . The value of the scalar triple product is the determinant with , and as its rows. With the -axis along the square of this determinant is
Repeating this calculation with the -axis along gives , while with the -axis along it is . Equating these expressions and dividing throughout by gives
where is the volume of the parallelepiped formed by the position vector of the vertices of the spherical triangle. Consequently, the result follows.
It is easy to see how for small spherical triangles, when the radius of the sphere is much greater than the sides of the triangle, this formula becomes the planar formula at the limit, since
and the same for and .
Geometric proof
Consider a unit sphere with:
Construct point and point such that
Construct point such that
It can therefore be seen that and
Notice that is the projection of on plane . Therefore
By basic trigonometry, we have:
But
Combining them we have:
By applying similar reasoning, we obtain the spherical law of sine:
Other proofs
A purely algebraic proof can be constructed from the spherical law of cosines. From the identity and the explicit expression for from the spherical law of cosines
Since the right hand side is invariant under a cyclic permutation of the spherical sine rule follows immediately.
The figure used in the Geometric proof above is used by and also provided in Banerjee (see Figure 3 in this paper) to derive the sine law using elementary linear algebra and projection matrices.
Hyperbolic case
In hyperbolic geometry when the curvature is −1, the law of sines becomes
In the special case when is a right angle, one gets
which is the analog of the formula in Euclidean geometry expressing the sine of an angle as the opposite side divided by the hypotenuse.
The case of surfaces of constant curvature
Define a generalized sine function, depending also on a real parameter :
The law of sines in constant curvature reads as
By substituting , , and , one obtains respectively , , and , that is, the Euclidean, spherical, and hyperbolic cases of the law of sines described above.
Let indicate the circumference of a circle of radius in a space of constant curvature . Then . Therefore, the law of sines can also be expressed as:
This formulation was discovered by János Bolyai.
Higher dimensions
A tetrahedron has four triangular facets. The absolute value of the polar sine () of the normal vectors to the three facets that share a vertex of the tetrahedron, divided by the area of the fourth facet will not depend upon the choice of the vertex:
More generally, for an -dimensional simplex (i.e., triangle (), tetrahedron (), pentatope (), etc.) in -dimensional Euclidean space, the absolute value of the polar sine of the normal vectors of the facets that meet at a vertex, divided by the hyperarea of the facet opposite the vertex is independent of the choice of the vertex. Writing for the hypervolume of the -dimensional simplex and for the product of the hyperareas of its -dimensional facets, the common ratio is
| Mathematics | Trigonometry | null |
113262 | https://en.wikipedia.org/wiki/Diprotodontia | Diprotodontia | Diprotodontia (, from Greek "two forward teeth") is the largest extant order of marsupials, with about 155 species, including the kangaroos, wallabies, possums, koala, wombats, and many others. Extinct diprotodonts include the hippopotamus-sized Diprotodon, and Thylacoleo, the so-called "marsupial lion".
Characteristics
Living diprotodonts are almost all herbivores, as were most of those that are now extinct. A few insectivorous and omnivorous diprotodonts are known, and the Potoroidae are almost unique among vertebrates in being largely fungivorous, but these seem to have arisen as relatively recent adaptations from the mainstream herbivorous lifestyle. The extinct thylacoleonids ("marsupial lions") are the only known group to have exhibited carnivory on a large scale.
Diprotodonts are restricted to Australasia. The earliest known fossils date to the late Oligocene, but their genesis certainly lies earlier than this, as large gaps occur in Australia's fossil record, with virtually no fossil record at all in geologically active New Guinea. The great diversity of known Oligocene diprotodonts suggests the order began to diverge well beforehand.
Many of the largest and least athletic diprotodonts (along with a wide range of other Australian megafauna) became extinct when humans first arrived in Australia about 50,000 years ago. Their extinction possibly occurred as a direct result of hunting, but was more probably a result of widespread habitat changes brought about by human activities—notably the use of fire.
Two key anatomical features, in combination, identify Diprotodontia. Members of the order are, first, "diprotodont" (meaning "two front teeth"): they have a pair of large, procumbent incisors on the lower jaw, a common feature of many early groups of mammals and mammaliforms. The diprotodont jaw is short, usually with three pairs of upper incisors (wombats, like rodents have only one pair), and no lower canines. The second trait distinguishing diprotodonts is "syndactyly", a fusing of the second and third digits of the foot up to the base of the claws, which leaves the claws themselves separate. Digit five is usually absent, and digit four is often greatly enlarged.
Syndactyly is not particularly common (though the Australian omnivorous marsupials share it) and is generally posited as an adaptation to assist in climbing. Many modern diprotodonts, however, are strictly terrestrial, and have evolved further adaptations to their feet to better suit this lifestyle. This makes the history of the tree-kangaroos particularly convoluted: it appears that the animals were arboreal at some time in the far distant past, moving afterward to the ground—gaining long kangaroo-like feet in the process — before returning to the trees, where they further developed a shortening and broadening of the hind feet and a novel climbing method.
Fossil record
The earliest known fossil of Diprotodontia dates back to the Late Oligocene (23.03 - 28.4 million years ago), and the earliest identifiable species is Hypsiprymnodon bartholomaii from the Early Miocene.
Classification
Until recently, only two suborders in Diprotodontia were noted: Vombatiformes which encompassed the wombats and koala and Phalangerida which contained all other families. Kirsch et al. (1997) split the families into three suborders. In addition, the six Phalangeriformes families are split into two superfamilies. The Macropodiformes are probably nested within the Phalangeriformes, though whether they are sister to Phalangeroidea or Petauroidea is debated.
Order Diprotodontia
Suborder Vombatiformes
Family Vombatidae: wombats (three species)
Family Phascolarctidae: koala (one species)
Family †Ilariidae
Family †Maradidae
Family †Diprotodontidae: (giant wombats)
Family †Palorchestidae: (marsupial tapirs)
Family †Thylacoleonidae: (marsupial lions)
Family †Wynyardiidae
Suborder Phalangeriformes
Superfamily Phalangeroidea
Family Phalangeridae: (brushtail possums and cuscuses)
Family Burramyidae: (pygmy possums)
Family †Ektopodontidae: (sprite possums)
Superfamily Petauroidea
Family Tarsipedidae: (honey possum)
Family Petauridae: (striped possum, Leadbeater's possum, yellow-bellied glider, sugar glider, mahogany glider, squirrel glider)
Family Pseudocheiridae: (ring-tailed possums and allies)
Family Acrobatidae: (feathertail glider and feather-tailed possum)
Suborder Macropodiformes
Family †Balbaridae: (basal quadrupedal kangaroos)
Family Macropodidae: (kangaroos, wallabies and allies)
Family Potoroidae: (bettongs, potoroos, and rat-kangaroos)
Family Hypsiprymnodontidae: (musky rat-kangaroo)
† means extinct family, genus or species
| Biology and health sciences | Marsupials | null |
113276 | https://en.wikipedia.org/wiki/Adze | Adze | An adze () or adz is an ancient and versatile cutting tool similar to an axe but with the cutting edge perpendicular to the handle rather than parallel. Adzes have been used since the Stone Age. They are used for smoothing or carving wood in hand woodworking, and as a hoe for agriculture and horticulture. Two basic forms of an adze are the hand adze (short hoe)—a short-handled tool swung with one hand—and the foot adze (hoe)—a long-handled tool capable of powerful swings using both hands, the cutting edge usually striking at foot or shin level. A similar tool is called a mattock, which differs by having two blades, one perpendicular to the handle and one parallel.
History
Africa
The adze is depicted in ancient Egyptian art from the Old Kingdom onward. Originally the adze blades were made of stone, but already in the Predynastic Period copper adzes had all but replaced those made of flint. Stone blades were fastened to the handle by tying and early bronze blades continued this simple construction. It was not until the later Bronze Age that the handle passes through an eye at the top of the blade. Examples of Egyptian adzes can be found in museums and on the Petrie Museum website.
A depiction of an adze was also used as a hieroglyph, representing the consonants stp, "chosen", and used as: ...Pharaoh XX, chosen of God/Goddess YY...
The ahnetjer (Manuel de Codage transliteration: aH-nTr) depicted as an adze-like instrument, was used in the Opening of the Mouth ceremony, intended to convey power over their senses to statues and mummies. It was apparently the foreleg of a freshly sacrificed bull or cow with which the mouth was touched.
As Iron Age technology moved south into Africa with migrating ancient Egyptians, they carried their technology with them, including adzes. To this day, iron adzes are used all over rural Africa for various purposes—from digging pit latrines, and chopping firewood, to tilling crop fields—whether they are of maize (corn), coffee, tea, pyrethrum, beans, millet, yams, or a plethora of other cash and subsistence crops.
New Zealand
Prehistoric Māori adzes from New Zealand were for wood carving, typically made from pounamu sourced from the South Island. During the Māori Archaic period found on the North Island were commonly made from greywacke from Motutapu Island or basalt from Opito Bay in the Coromandel, similar to adzes constructed on other Pacific Islands. Early period notched adzes found in Northland were primarily made of argillite quarried from locations around the Marlborough and Nelson regions. At the same time on Henderson Island, a small coral island in eastern Polynesia lacking any rock other than limestone, native populations may have fashioned giant clamshells into adzes.
Northwest Coastal America
American Northwest coast native peoples traditionally used adzes for both functional construction (from bowls to canoes) and art (from masks to totem poles). Northwest coast adzes take two forms: hafted and D-handle. The hafted form is similar in form to a European adze with the haft constructed from a natural crooked branch which approximately forms a 60% angle. The thin end is used as the handle and the thick end is flattened and notched such that an adze iron can be lashed to it. Modern hafts are sometimes constructed from a sawed blank with a dowel added for strength at the crook. The second form is the D-handle adze which is basically an adze iron with a directly attached handle. The D-handle, therefore, provides no mechanical leverage. Northwest coast adzes are often classified by size and iron shape vs. role. As with European adzes, iron shapes include straight, gutter and lipped. Where larger Northwest adzes are similar in size to their European counterparts, the smaller sizes are typically much lighter such that they can be used for the detailed smoothing, shaping and surface texturing required for figure carving. Final surfacing is sometimes performed with a crooked knife.
New Guinea and Melanesia
Ground stone adzes used to be produced by a variety of people in Western New Guinea (Indonesia), Papua New Guinea and some of the smaller Islands of Melanesia and Micronesia. The hardstone would have been ground on a riverine rock with the help of water until the desired shape was obtained. It was then fixed to a natural grown angled wood with resin and plant fibers. A variety of minerals were used. Imported steel axes or machetes have now entirely replaced these tools for decades in even the remotest parts of New Guinea. Indeed, even before the first foreign missionaries or colonial officials arrived in the New Guinea Highlands, inhabitants had already obtained steel tools through trade with their neighbors. Stone tools are sometimes manufactured to be sold as curios to tourists.
Modern adzes
Modern adzes are made from steel with wooden handles, and enjoy limited use: occasionally in semi-industrial areas, but particularly by "revivalists" such as those at the Colonial Williamsburg cultural center in Virginia, United States. However, the traditional adze has largely been replaced by the sawmill and the powered plane, at least in industrialised cultures. It remains in use for some specialist crafts, for example by coopers. Adzes are also in current use by artists such as Northwest Coast American and Canadian Indigenous sculptors doing totem pole carving, as well as masks and bowls.
Foot adze
"Adzes are used for removing heavy waste, leveling, shaping, or trimming the surfaces of timber" and boards. Generally, the user stands astride a board or log and swings the adze downwards between his feet, chipping off pieces of wood, moving backwards as they go and leaving a relatively smooth surface behind.
Foot adzes are most commonly known as shipbuilder's or carpenter's adzes. They range in size from 00 to 5 being with the cutting edge wide. On the modern, steel adze the cutting edge may be flat for smoothing work to very rounded for hollowing work such as bowls, gutters and canoes. The shoulders or sides of an adze may be curved called a lipped adze, used for notching. The end away from the cutting edge is called the pole and be of different shapes, generally flat or a pin pole.
Carpenter's adze A heavy adze, often with very steep curves, and a very heavy, blunt pole. The weight of this adze makes it unsuitable for sustained overhead adzing.
Railroad adze A carpenter's adze which had its bit extended in an effort to limit the breaking of handles when shaping railroad ties (railway sleepers). Early examples in New England began showing up approximately in the 1840s–1850s. The initial prototypes clearly showed a weld where the extension was attached.
Shipwright's adze A lighter, and more versatile adze than the carpenter's adze. This was designed to be used in a variety of positions, including overhead, as well as in front on waist and chest level.
Lipped shipwright's adze A variation of the shipwright's adze. It features a wider than normal bit, whose outside edges are sharply turned up, so that when gazing directly down the adze, from bit to eye, the cutting edge resembles an extremely wide and often very flat U. This adze was mainly used for shaping cross grain, such as for joining planks.
Another group of adzes can be differentiated by the handles; the D-handled adzes have a handle where the hand can be wrapped around the D, close to the bit. These adzes closely follow traditional forms in that the bit or tooth is not wrapped around the handle as a head.
The head of an ice axe typically possesses an adze for chopping rough steps in ice.
A firefighter tool called the Halligan bar has a dull adze on one end of the bar. This bar is a multipurpose tool for forcible entry of a structure and demolition with a forked pry-bar on one end and an adze and spike on the other, called the adze-end.
Demolition adze A demolition adze has a dull edge and is used for separating materials in the demolition or salvage of old buildings.
Hand adze
There are a number of specialist, short-handled adzes used by coopers, wainwrights, and chair makers, and in bowl and trough making. Many of these have shorter handles for control and more curve in the head to allow better clearance for shorter cuts.
Bulgarian adze
During the communist period of Bulgaria, a new multi-use woodworking adze, called (), emerged. It has a sharpened edge perpendicular to the handle, resembling an adze, but it is also used like a carpenter's hammer. On the back of the head is a textured poll for driving nails, and on the front is a V-shaped hole used for prying, to extract the bent nails.
Some urban legends say that Bulgarian migrant workers always carry their adzes with them so they can do construction work more efficiently due to the lack of Western equivalent of the tool. The Bulgarian adze is often mistaken for a hammer.
There is a popular Bulgarian folk song called "На теслата дръжката" (eng: The tesla's handle) about a craftsman and the masculinity of his tool.
| Technology | Hand tools | null |
113302 | https://en.wikipedia.org/wiki/Surface%20tension | Surface tension | Surface tension is the tendency of liquid surfaces at rest to shrink into the minimum surface area possible. Surface tension is what allows objects with a higher density than water such as razor blades and insects (e.g. water striders) to float on a water surface without becoming even partly submerged.
At liquid–air interfaces, surface tension results from the greater attraction of liquid molecules to each other (due to cohesion) than to the molecules in the air (due to adhesion).
There are two primary mechanisms in play. One is an inward force on the surface molecules causing the liquid to contract. Second is a tangential force parallel to the surface of the liquid. This tangential force is generally referred to as the surface tension. The net effect is the liquid behaves as if its surface were covered with a stretched elastic membrane. But this analogy must not be taken too far as the tension in an elastic membrane is dependent on the amount of deformation of the membrane while surface tension is an inherent property of the liquid–air or liquid–vapour interface.
Because of the relatively high attraction of water molecules to each other through a web of hydrogen bonds, water has a higher surface tension (72.8 millinewtons (mN) per meter at 20 °C) than most other liquids. Surface tension is an important factor in the phenomenon of capillarity.
Surface tension has the dimension of force per unit length, or of energy per unit area. The two are equivalent, but when referring to energy per unit of area, it is common to use the term surface energy, which is a more general term in the sense that it applies also to solids.
In materials science, surface tension is used for either surface stress or surface energy.
Causes
Due to the cohesive forces, a molecule located away from the surface is pulled equally in every direction by neighboring liquid molecules, resulting in a net force of zero. The molecules at the surface do not have the same molecules on all sides of them and therefore are pulled inward. This creates some internal pressure and forces liquid surfaces to contract to the minimum area.
There is also a tension parallel to the surface at the liquid-air interface which will resist an external force, due to the cohesive nature of water molecules.
The forces of attraction acting between molecules of the same type are called cohesive forces, while those acting between molecules of different types are called adhesive forces. The balance between the cohesion of the liquid and its adhesion to the material of the container determines the degree of wetting, the contact angle, and the shape of meniscus. When cohesion dominates (specifically, adhesion energy is less than half of cohesion energy) the wetting is low and the meniscus is convex at a vertical wall (as for mercury in a glass container). On the other hand, when adhesion dominates (when adhesion energy is more than half of cohesion energy) the wetting is high and the similar meniscus is concave (as in water in a glass).
Surface tension is responsible for the shape of liquid droplets. Although easily deformed, droplets of water tend to be pulled into a spherical shape by the imbalance in cohesive forces of the surface layer. In the absence of other forces, drops of virtually all liquids would be approximately spherical. The spherical shape minimizes the necessary "wall tension" of the surface layer according to Laplace's law.
Another way to view surface tension is in terms of energy. A molecule in contact with a neighbor is in a lower state of energy than if it were alone. The interior molecules have as many neighbors as they can possibly have, but the boundary molecules are missing neighbors (compared to interior molecules) and therefore have higher energy. For the liquid to minimize its energy state, the number of higher energy boundary molecules must be minimized. The minimized number of boundary molecules results in a minimal surface area.
As a result of surface area minimization, a surface will assume a smooth shape.
Physics
Physical units
Surface tension, represented by the symbol γ (alternatively σ or T), is measured in force per unit length. Its SI unit is newton per meter but the cgs unit of dyne per centimeter is also used. For example,
Definition
Surface tension can be defined in terms of force or energy.
In terms of force
Surface tension of a liquid is the force per unit length. In the illustration on the right, the rectangular frame, composed of three unmovable sides (black) that form a "U" shape, and a fourth movable side (blue) that can slide to the right. Surface tension will pull the blue bar to the left; the force required to hold the movable side is proportional to the length of the immobile side. Thus the ratio depends only on the intrinsic properties of the liquid (composition, temperature, etc.), not on its geometry. For example, if the frame had a more complicated shape, the ratio , with the length of the movable side and the force required to stop it from sliding, is found to be the same for all shapes. We therefore define the surface tension as
The reason for the is that the film has two sides (two surfaces), each of which contributes equally to the force; so the force contributed by a single side is .
In terms of energy
Surface tension of a liquid is the ratio of the change in the energy of the liquid to the change in the surface area of the liquid (that led to the change in energy). This can be easily related to the previous definition in terms of force: if is the force required to stop the side from starting to slide, then this is also the force that would keep the side in the state of sliding at a constant speed (by Newton's Second Law). But if the side is moving to the right (in the direction the force is applied), then the surface area of the stretched liquid is increasing while the applied force is doing work on the liquid. This means that increasing the surface area increases the energy of the film. The work done by the force in moving the side by distance is ; at the same time the total area of the film increases by (the factor of 2 is here because the liquid has two sides, two surfaces). Thus, multiplying both the numerator and the denominator of by , we get
This work is, by the usual arguments, interpreted as being stored as potential energy. Consequently, surface tension can be also measured in SI system as joules per square meter and in the cgs system as ergs per cm2. Since mechanical systems try to find a state of minimum potential energy, a free droplet of liquid naturally assumes a spherical shape, which has the minimum surface area for a given volume.
The equivalence of measurement of energy per unit area to force per unit length can be proven by dimensional analysis.
Effects
Water
Several effects of surface tension can be seen with ordinary water:
Surfactants
Surface tension is visible in other common phenomena, especially when surfactants are used to decrease it:
Soap bubbles have very large surface areas with very little mass. Bubbles in pure water are unstable. The addition of surfactants, however, can have a stabilizing effect on the bubbles (see Marangoni effect). Surfactants actually reduce the surface tension of water by a factor of three or more.
Emulsions are a type of colloidal dispersion in which surface tension plays a role. Tiny droplets of oil dispersed in pure water will spontaneously coalesce and phase separate. The addition of surfactants reduces the interfacial tension and allow for the formation of oil droplets in the water medium (or vice versa). The stability of such formed oil droplets depends on many different chemical and environmental factors.
Surface curvature and pressure
If no force acts normal to a tensioned surface, the surface must remain flat. But if the pressure on one side of the surface differs from pressure on the other side, the pressure difference times surface area results in a normal force. In order for the surface tension forces to cancel the force due to pressure, the surface must be curved. The diagram shows how surface curvature of a tiny patch of surface leads to a net component of surface tension forces acting normal to the center of the patch. When all the forces are balanced, the resulting equation is known as the Young–Laplace equation:
where:
is the pressure difference, known as the Laplace pressure.
is surface tension.
and are radii of curvature in each of the axes that are parallel to the surface.
The quantity in parentheses on the right hand side is in fact (twice) the mean curvature of the surface (depending on normalisation).
Solutions to this equation determine the shape of water drops, puddles, menisci, soap bubbles, and all other shapes determined by surface tension (such as the shape of the impressions that a water strider's feet make on the surface of a pond).
The table below shows how the internal pressure of a water droplet increases with decreasing radius. For not very small drops the effect is subtle, but the pressure difference becomes enormous when the drop sizes approach the molecular size. (In the limit of a single molecule the concept becomes meaningless.)
Floating objects
When an object is placed on a liquid, its weight depresses the surface, and if surface tension and downward force become equal then it is balanced by the surface tension forces on either side , which are each parallel to the water's surface at the points where it contacts the object. Notice that small movement in the body may cause the object to sink. As the angle of contact decreases, surface tension decreases. The horizontal components of the two arrows point in opposite directions, so they cancel each other, but the vertical components point in the same direction and therefore add up to balance . The object's surface must not be wettable for this to happen, and its weight must be low enough for the surface tension to support it. If denotes the mass of the needle and acceleration due to gravity, we have
Liquid surface
To find the shape of the minimal surface bounded by some arbitrary shaped frame using strictly mathematical means can be a daunting task. Yet by fashioning the frame out of wire and dipping it in soap-solution, a locally minimal surface will appear in the resulting soap-film within seconds.
The reason for this is that the pressure difference across a fluid interface is proportional to the mean curvature, as seen in the Young–Laplace equation. For an open soap film, the pressure difference is zero, hence the mean curvature is zero, and minimal surfaces have the property of zero mean curvature.
Contact angles
The surface of any liquid is an interface between that liquid and some other medium. The top surface of a pond, for example, is an interface between the pond water and the air. Surface tension, then, is not a property of the liquid alone, but a property of the liquid's interface with another medium. If a liquid is in a container, then besides the liquid/air interface at its top surface, there is also an interface between the liquid and the walls of the container. The surface tension between the liquid and air is usually different (greater) than its surface tension with the walls of a container. And where the two surfaces meet, their geometry must be such that all forces balance.
Where the two surfaces meet, they form a contact angle, , which is the angle the tangent to the surface makes with the solid surface. Note that the angle is measured through the liquid, as shown in the diagrams above. The diagram to the right shows two examples. Tension forces are shown for the liquid–air interface, the liquid–solid interface, and the solid–air interface. The example on the left is where the difference between the liquid–solid and solid–air surface tension, , is less than the liquid–air surface tension, , but is nevertheless positive, that is
In the diagram, both the vertical and horizontal forces must cancel exactly at the contact point, known as equilibrium. The horizontal component of is canceled by the adhesive force, .
The more telling balance of forces, though, is in the vertical direction. The vertical component of must exactly cancel the difference of the forces along the solid surface, .
Since the forces are in direct proportion to their respective surface tensions, we also have:
where
is the liquid–solid surface tension,
is the liquid–air surface tension,
is the solid–air surface tension,
is the contact angle, where a concave meniscus has contact angle less than 90° and a convex meniscus has contact angle of greater than 90°.
This means that although the difference between the liquid–solid and solid–air surface tension, , is difficult to measure directly, it can be inferred from the liquid–air surface tension, , and the equilibrium contact angle, , which is a function of the easily measurable advancing and receding contact angles (see main article contact angle).
This same relationship exists in the diagram on the right. But in this case we see that because the contact angle is less than 90°, the liquid–solid/solid–air surface tension difference must be negative:
Special contact angles
Observe that in the special case of a water–silver interface where the contact angle is equal to 90°, the liquid–solid/solid–air surface tension difference is exactly zero.
Another special case is where the contact angle is exactly 180°. Water with specially prepared Teflon approaches this. Contact angle of 180° occurs when the liquid–solid surface tension is exactly equal to the liquid–air surface tension.
Liquid in a vertical tube
An old style mercury barometer consists of a vertical glass tube about 1 cm in diameter partially filled with mercury, and with a vacuum (called Torricelli's vacuum) in the unfilled volume (see diagram to the right). Notice that the mercury level at the center of the tube is higher than at the edges, making the upper surface of the mercury dome-shaped. The center of mass of the entire column of mercury would be slightly lower if the top surface of the mercury were flat over the entire cross-section of the tube. But the dome-shaped top gives slightly less surface area to the entire mass of mercury. Again the two effects combine to minimize the total potential energy. Such a surface shape is known as a convex meniscus.
We consider the surface area of the entire mass of mercury, including the part of the surface that is in contact with the glass, because mercury does not adhere to glass at all. So the surface tension of the mercury acts over its entire surface area, including where it is in contact with the glass. If instead of glass, the tube was made out of copper, the situation would be very different. Mercury aggressively adheres to copper. So in a copper tube, the level of mercury at the center of the tube will be lower than at the edges (that is, it would be a concave meniscus). In a situation where the liquid adheres to the walls of its container, we consider the part of the fluid's surface area that is in contact with the container to have negative surface tension. The fluid then works to maximize the contact surface area. So in this case increasing the area in contact with the container decreases rather than increases the potential energy. That decrease is enough to compensate for the increased potential energy associated with lifting the fluid near the walls of the container.
If a tube is sufficiently narrow and the liquid adhesion to its walls is sufficiently strong, surface tension can draw liquid up the tube in a phenomenon known as capillary action. The height to which the column is lifted is given by Jurin's law:
where
is the height the liquid is lifted,
is the liquid–air surface tension,
is the density of the liquid,
is the radius of the capillary,
is the acceleration due to gravity,
is the angle of contact described above. If is greater than 90°, as with mercury in a glass container, the liquid will be depressed rather than lifted.
Puddles on a surface
Pouring mercury onto a horizontal flat sheet of glass results in a puddle that has a perceptible thickness. The puddle will spread out only to the point where it is a little under half a centimetre thick, and no thinner. Again this is due to the action of mercury's strong surface tension. The liquid mass flattens out because that brings as much of the mercury to as low a level as possible, but the surface tension, at the same time, is acting to reduce the total surface area. The result of the compromise is a puddle of a nearly fixed thickness.
The same surface tension demonstration can be done with water, lime water or even saline, but only on a surface made of a substance to which water does not adhere. Wax is such a substance. Water poured onto a smooth, flat, horizontal wax surface, say a waxed sheet of glass, will behave similarly to the mercury poured onto glass.
The thickness of a puddle of liquid on a surface whose contact angle is 180° is given by:
where
is the depth of the puddle in centimeters or meters.
is the surface tension of the liquid in dynes per centimeter or newtons per meter.
is the acceleration due to gravity and is equal to 980 cm/s2 or 9.8 m/s2
is the density of the liquid in grams per cubic centimeter or kilograms per cubic meter
In reality, the thicknesses of the puddles will be slightly less than what is predicted by the above formula because very few surfaces have a contact angle of 180° with any liquid. When the contact angle is less than 180°, the thickness is given by:
For mercury on glass, = 487 dyn/cm, = 13.5 g/cm3 and = 140°, which gives = 0.36 cm. For water on paraffin at 25 °C, = 72 dyn/cm, = 1.0 g/cm3, and = 107° which gives = 0.44 cm.
The formula also predicts that when the contact angle is 0°, the liquid will spread out into a micro-thin layer over the surface. Such a surface is said to be fully wettable by the liquid.
Breakup of streams into drops
In day-to-day life all of us observe that a stream of water emerging from a faucet will break up into droplets, no matter how smoothly the stream is emitted from the faucet. This is due to a phenomenon called the Plateau–Rayleigh instability, which is entirely a consequence of the effects of surface tension.
The explanation of this instability begins with the existence of tiny perturbations in the stream. These are always present, no matter how smooth the stream is. If the perturbations are resolved into sinusoidal components, we find that some components grow with time while others decay with time. Among those that grow with time, some grow at faster rates than others. Whether a component decays or grows, and how fast it grows is entirely a function of its wave number (a measure of how many peaks and troughs per centimeter) and the radii of the original cylindrical stream.
Gallery
Thermodynamics
Thermodynamic theories of surface tension
J.W. Gibbs developed the thermodynamic theory of capillarity based
on the idea of surfaces of discontinuity. Gibbs considered the case of a sharp mathematical surface being placed somewhere within the microscopically fuzzy physical interface that exists between two homogeneous substances. Realizing that the exact choice of the surface's location was somewhat arbitrary, he left it flexible. Since the interface exists in thermal and chemical equilibrium with the substances around it (having temperature and chemical potentials ), Gibbs considered the case where the surface may have excess energy, excess entropy, and excess particles, finding the natural free energy function in this case to be , a quantity later named as the grand potential and given the symbol .
Considering a given subvolume containing a surface of discontinuity, the volume is divided by the mathematical surface into two parts A and B, with volumes and , with exactly. Now, if the two parts A and B were homogeneous fluids (with pressures , ) and remained perfectly homogeneous right up to the mathematical boundary, without any surface effects, the total grand potential of this volume would be simply . The surface effects of interest are a modification to this, and they can be all collected into a surface free energy term so the total grand potential of the volume becomes:
For sufficiently macroscopic and gently curved surfaces, the surface free energy must simply be proportional to the surface area:
for surface tension and surface area .
As stated above, this implies the mechanical work needed to increase a surface area A is , assuming the volumes on each side do not change. Thermodynamics requires that for systems held at constant chemical potential and temperature, all spontaneous changes of state are accompanied by a decrease in this free energy , that is, an increase in total entropy taking into account the possible movement of energy and particles from the surface into the surrounding fluids. From this it is easy to understand why decreasing the surface area of a mass of liquid is always spontaneous, provided it is not coupled to any other energy changes. It follows that in order to increase surface area, a certain amount of energy must be added.
Gibbs and other scientists have wrestled with the arbitrariness in the exact microscopic placement of the surface. For microscopic surfaces with very tight curvatures, it is not correct to assume the surface tension is independent of size, and topics like the Tolman length come into play. For a macroscopic-sized surface (and planar surfaces), the surface placement does not have a significant effect on ; however, it does have a very strong effect on the values of the surface entropy, surface excess mass densities, and surface internal energy, which are the partial derivatives of the surface tension function .
Gibbs emphasized that for solids, the surface free energy may be completely different from surface stress (what he called surface tension): the surface free energy is the work required to form the surface, while surface stress is the work required to stretch the surface. In the case of a two-fluid interface, there is no distinction between forming and stretching because the fluids and the surface completely replenish their nature when the surface is stretched. For a solid, stretching the surface, even elastically, results in a fundamentally changed surface. Further, the surface stress on a solid is a directional quantity (a stress tensor) while surface energy is scalar.
Fifteen years after Gibbs, J.D. van der Waals developed the theory of capillarity effects based on the hypothesis of a continuous variation of density. He added to the energy density the term where c is the capillarity coefficient and ρ is the density. For the multiphase equilibria, the results of the van der Waals approach practically coincide with the Gibbs formulae, but for modelling of the dynamics of phase transitions the van der Waals approach is much more convenient. The van der Waals capillarity energy is now widely used in the phase field models of multiphase flows. Such terms are also discovered in the dynamics of non-equilibrium gases.
Thermodynamics of bubbles
The pressure inside an ideal spherical bubble can be derived from thermodynamic free energy considerations. The above free energy can be written as:
where is the pressure difference between the inside (A) and outside (B) of the bubble, and is the bubble volume. In equilibrium, , and so,
For a spherical bubble, the volume and surface area are given simply by
and
Substituting these relations into the previous expression, we find
which is equivalent to the Young–Laplace equation when .
Influence of temperature
Surface tension is dependent on temperature. For that reason, when a value is given for the surface tension of an interface, temperature must be explicitly stated. The general trend is that surface tension decreases with the increase of temperature, reaching a value of 0 at the critical temperature. For further details see Eötvös rule. There are only empirical equations to relate surface tension and temperature:
Eötvös: Here is the molar volume of a substance, is the critical temperature and is a constant valid for almost all substances. A typical value is = . For water one can further use = 18 ml/mol and = 647 K (374 °C). A variant on Eötvös is described by Ramay and Shields: where the temperature offset of 6 K provides the formula with a better fit to reality at lower temperatures.
Guggenheim–Katayama: is a constant for each liquid and is an empirical factor, whose value is for organic liquids. This equation was also proposed by van der Waals, who further proposed that could be given by the expression where is a universal constant for all liquids, and is the critical pressure of the liquid (although later experiments found to vary to some degree from one liquid to another).
Both Guggenheim–Katayama and Eötvös take into account the fact that surface tension reaches 0 at the critical temperature, whereas Ramay and Shields fails to match reality at this endpoint.
Influence of solute concentration
Solutes can have different effects on surface tension depending on the nature of the surface and the solute:
Little or no effect, for example sugar at water|air, most organic compounds at oil/air
Increase surface tension, most inorganic salts at water|air
Non-monotonic change, most inorganic acids at water|air
Decrease surface tension progressively, as with most amphiphiles, e.g., alcohols at water|air
Decrease surface tension until certain critical concentration, and no effect afterwards: surfactants that form micelles
What complicates the effect is that a solute can exist in a different concentration at the surface of a solvent than in its bulk. This difference varies from one solute–solvent combination to another.
Gibbs isotherm states that:
is known as surface concentration, it represents excess of solute per unit area of the surface over what would be present if the bulk concentration prevailed all the way to the surface. It has units of mol/m2
is the concentration of the substance in the bulk solution.
is the gas constant and the temperature
Certain assumptions are taken in its deduction, therefore Gibbs isotherm can only be applied to ideal (very dilute) solutions with two components.
Influence of particle size on vapor pressure
The Clausius–Clapeyron relation leads to another equation also attributed to Kelvin, as the Kelvin equation. It explains why, because of surface tension, the vapor pressure for small droplets of liquid in suspension is greater than standard vapor pressure of that same liquid when the interface is flat. That is to say that when a liquid is forming small droplets, the equilibrium concentration of its vapor in its surroundings is greater. This arises because the pressure inside the droplet is greater than outside.
is the standard vapor pressure for that liquid at that temperature and pressure.
is the molar volume.
is the gas constant
is the Kelvin radius, the radius of the droplets.
The effect explains supersaturation of vapors. In the absence of nucleation sites, tiny droplets must form before they can evolve into larger droplets. This requires a vapor pressure many times the vapor pressure at the phase transition point.
This equation is also used in catalyst chemistry to assess mesoporosity for solids.
The effect can be viewed in terms of the average number of molecular neighbors of surface molecules (see diagram).
The table shows some calculated values of this effect for water at different drop sizes:
The effect becomes clear for very small drop sizes, as a drop of 1 nm radius has about 100 molecules inside, which is a quantity small enough to require a quantum mechanics analysis.
Methods of measurement
Because surface tension manifests itself in various effects, it offers a number of paths to its measurement. Which method is optimal depends upon the nature of the liquid being measured, the conditions under which its tension is to be measured, and the stability of its surface when it is deformed. An instrument that measures surface tension is called tensiometer.
Du Noüy ring method: The traditional method used to measure surface or interfacial tension. Wetting properties of the surface or interface have little influence on this measuring technique. Maximum pull exerted on the ring by the surface is measured.
Wilhelmy plate method: A universal method especially suited to check surface tension over long time intervals. A vertical plate of known perimeter is attached to a balance, and the force due to wetting is measured.
Spinning drop method: This technique is ideal for measuring low interfacial tensions. The diameter of a drop within a heavy phase is measured while both are rotated.
Pendant drop method: Surface and interfacial tension can be measured by this technique, even at elevated temperatures and pressures. Geometry of a drop is analyzed optically. For pendant drops the maximum diameter and the ratio between this parameter and the diameter at the distance of the maximum diameter from the drop apex has been used to evaluate the size and shape parameters in order to determine surface tension.
Bubble pressure method (Jaeger's method): A measurement technique for determining surface tension at short surface ages. Maximum pressure of each bubble is measured.
Drop volume method: A method for determining interfacial tension as a function of interface age. Liquid of one density is pumped into a second liquid of a different density and time between drops produced is measured.
Capillary rise method: The end of a capillary is immersed into the solution. The height at which the solution reaches inside the capillary is related to the surface tension by the equation discussed above.
Stalagmometric method: A method of weighting and reading a drop of liquid.
Sessile drop method: A method for determining surface tension and density by placing a drop on a substrate and measuring the contact angle (see Sessile drop technique).
Du Noüy–Padday method: A minimized version of Du Noüy method uses a small diameter metal needle instead of a ring, in combination with a high sensitivity microbalance to record maximum pull. The advantage of this method is that very small sample volumes (down to few tens of microliters) can be measured with very high precision, without the need to correct for buoyancy (for a needle or rather, rod, with proper geometry). Further, the measurement can be performed very quickly, minimally in about 20 seconds.
Vibrational frequency of levitated drops: The natural frequency of vibrational oscillations of magnetically levitated drops has been used to measure the surface tension of superfluid 4He. This value is estimated to be 0.375 dyn/cm at = 0 K.
Resonant oscillations of spherical and hemispherical liquid drop: The technique is based on measuring the resonant frequency of spherical and hemispherical pendant droplets driven in oscillations by a modulated electric field. The surface tension and viscosity can be evaluated from the obtained resonant curves.
Drop-bounce method: This method is based on aerodynamic levitation with a split-able nozzle design. After dropping a stably levitated droplet onto a platform, the sample deforms and bounces back, oscillating in mid-air as it tries to minimize its surface area. Through this oscillation behavior, the liquid's surface tension and viscosity can be measured.
Values
Data table
Surface tension of water
The surface tension of pure liquid water in contact with its vapor has been given by IAPWS as
where both and the critical temperature = 647.096 K are expressed in kelvins. The region of validity the entire vapor–liquid saturation curve, from the triple point (0.01 °C) to the critical point. It also provides reasonable results when extrapolated to metastable (supercooled) conditions, down to at least −25 °C. This formulation was originally adopted by IAPWS in 1976 and was adjusted in 1994 to conform to the International Temperature Scale of 1990.
The uncertainty of this formulation is given over the full range of temperature by IAPWS. For temperatures below 100 °C, the uncertainty is ±0.5%.
Surface tension of seawater
Nayar et al. published reference data for the surface tension of seawater over the salinity range of and a temperature range of at atmospheric pressure. The range of temperature and salinity encompasses both the oceanographic range and the range of conditions encountered in thermal desalination technologies. The uncertainty of the measurements varied from 0.18 to 0.37 mN/m with the average uncertainty being 0.22 mN/m.
Nayar et al. correlated the data with the following equation
where is the surface tension of seawater in mN/m, is the surface tension of water in mN/m, is the reference salinity in g/kg, and is temperature in degrees Celsius. The average absolute percentage deviation between measurements and the correlation was 0.19% while the maximum deviation is 0.60%.
The International Association for the Properties of Water and Steam (IAPWS) has adopted this correlation as an international standard guideline.
| Physical sciences | Fluid mechanics | null |
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